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from diffusers import StableDiffusionXLPipeline, DDIMScheduler | |
import torch | |
import gradio as gr | |
import inversion | |
import numpy as np | |
from PIL import Image | |
import sa_handler | |
device = "cuda" if torch.cuda.is_available() else "cpu" | |
scheduler = DDIMScheduler(beta_start=0.00085, beta_end=0.012, beta_schedule="scaled_linear", clip_sample=False, set_alpha_to_one=False) | |
pipeline = StableDiffusionXLPipeline.from_pretrained("stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.float16, variant="fp16", use_safetensors=True, scheduler=scheduler).to(device) | |
def run(image, src_style, src_prompt, prompts, shared_score_shift, shared_score_scale, guidance_scale, num_inference_steps, large, seed): | |
prompts = prompts.splitlines() | |
dim, d = (1024, 128) if large else (512, 64) | |
image = image.resize((dim, dim)) | |
x0 = np.array(image) | |
zts = inversion.ddim_inversion(pipeline, x0, src_prompt, num_inference_steps, 2) | |
offset = min(5, len(zts) - 1) | |
prompts.insert(0, src_prompt) | |
shared_score_shift = np.log(shared_score_shift) | |
handler = sa_handler.Handler(pipeline) | |
sa_args = sa_handler.StyleAlignedArgs( | |
share_group_norm=True, share_layer_norm=True, share_attention=True, | |
adain_queries=True, adain_keys=True, adain_values=False, | |
shared_score_shift=shared_score_shift, shared_score_scale=shared_score_scale,) | |
handler.register(sa_args) | |
for i in range(1, len(prompts)): | |
prompts[i] = f'{prompts[i]}, {src_style}.' | |
zT, inversion_callback = inversion.make_inversion_callback(zts, offset=offset) | |
g_cpu = torch.Generator(device='cpu') | |
if seed > 0: | |
g_cpu.manual_seed(seed) | |
latents = torch.randn(len(prompts), 4, d, d, device='cpu', generator=g_cpu, dtype=pipeline.unet.dtype,).to(device) | |
latents[0] = zT | |
images_a = pipeline(prompts, latents=latents, callback_on_step_end=inversion_callback, num_inference_steps=num_inference_steps, guidance_scale=guidance_scale).images | |
handler.remove() | |
torch.cuda.empty_cache() | |
images_pil = [Image.fromarray((img * 255).astype(np.uint8)) for img in images_a] | |
return images_pil | |
with gr.Blocks() as demo: | |
gr.Markdown('''# Welcome to Tonic's Stable Style Align | |
Here you can generate images with a style from a reference image using [transfer style from sdxl](https://huggingface.co/docs/diffusers/main/en/using-diffusers/sdxl). Add a reference picture, describe the style and add prompts to generate images in that style. It's the most interesting with your own art!''') | |
with gr.Row(): | |
image_input = gr.Image(label="Reference image", type="pil") | |
with gr.Row(): | |
style_input = gr.Textbox(label="Describe the reference style") | |
image_desc_input = gr.Textbox(label="Describe the reference image") | |
prompts_input = gr.Textbox(label="Prompts to generate images (separate with new lines)", lines=5) | |
with gr.Accordion(label="Advanced Settings"): | |
with gr.Row(): | |
shared_score_shift_input = gr.Slider(value=1.1, label="shared_score_shift", minimum=1.0, maximum=2.0, step=0.05) | |
shared_score_scale_input = gr.Slider(value=1.0, label="shared_score_scale", minimum=0.0, maximum=1.0, step=0.05) | |
guidance_scale_input = gr.Slider(value=10.0, label="guidance_scale", minimum=5.0, maximum=20.0, step=1) | |
num_inference_steps_input = gr.Slider(value=12, label="num_inference_steps", minimum=1, maximum=12, step=1) | |
large_input = gr.Checkbox(False, label="Large (1024x1024)") | |
seed_input = gr.Slider(value=0, label="seed (0 for random)", minimum=0, maximum=1000000, step=42) | |
with gr.Row(): | |
run_button = gr.Button("Generate Images") | |
with gr.Row(): | |
output_gallery = gr.Gallery() | |
run_button.click( | |
run, | |
inputs=[image_input, style_input, image_desc_input, prompts_input, shared_score_shift_input, shared_score_scale_input, guidance_scale_input, num_inference_steps_input, large_input, seed_input], | |
outputs=output_gallery | |
) | |
demo.launch() |