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--- |
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license: other |
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base_model: diffusers/stable-diffusion-xl-base-0.9 |
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tags: |
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- stable-diffusion-xl |
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- stable-diffusion-xl-diffusers |
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- stable-diffusion |
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- text-to-image |
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- diffusers |
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- di.ffusion.ai |
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inference: true |
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extra_gated_prompt: >- |
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Copyright (c) Stability AI Ltd. and Source Code Bulgaria Ltd. |
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12.ADDITIONAL TERMS FOR FFUSION.AI |
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In addition to the terms set forth above, the following terms apply specifically |
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a. Any use of the FFusionXL-09-SDXL model must include proper attribution to |
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--- |
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# FFXL Model Card |
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<div style="display: flex; flex-wrap: wrap; gap: 2px;"> |
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<img src="https://img.shields.io/badge/%F0%9F%94%A5%20Refiner%20Compatible-Yes-success"> |
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<img src="https://img.shields.io/badge/%F0%9F%92%BB%20CLIP--ViT%2FG%20and%20CLIP--ViT%2FL%20tested-Yes-success"> |
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<img src="https://img.shields.io/badge/%F0%9F%A7%A8%20FFXL%20Diffusers-available-brightgreen"> |
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</div> |
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 |
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## Model |
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FFXL based on SDXL consists of a two-step pipeline for latent diffusion: |
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First, we use a base model to generate latents of the desired output size. |
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In the second step, we use a specialized high-resolution model and apply a technique called SDEdit (https://arxiv.org/abs/2108.01073, also known as "img2img") |
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to the latents generated in the first step, using the same prompt. |
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[](https://huggingface.co/FFusion/FFusionXL-09-SDXL/blob/main/FFusionXL-09-SDXL.safetensors) |
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### Model Description |
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- **Trained by:** FFusion AI |
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- **Model type:** Diffusion-based text-to-image generative model |
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- **License:** [FFXL Research License](https://huggingface.co/FFusion/FFusionXL-09-SDXL/blob/main/LICENSE.md) |
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- **Model Description:** This is a trained model based on SDXL that can be used to generate and modify images based on text prompts. It is a [Latent Diffusion Model](https://arxiv.org/abs/2112.10752) that uses two fixed, pretrained text encoders ([OpenCLIP-ViT/G](https://github.com/mlfoundations/open_clip) and [CLIP-ViT/L](https://github.com/openai/CLIP/tree/main)). |
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- **Resources for more information:** [SDXL paper on arXiv](https://arxiv.org/abs/2307.01952). |
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 |
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### Model Sources |
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- **Demo:** [](https://huggingface.co/spaces/FFusion/FFusionXL-SDXL-DEMO) |
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<div style="display: flex; flex-wrap: wrap; gap: 2px;"> |
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<a href="https://huggingface.co/FFusion/FFusion-BaSE" target="_new" rel="ugc"><img src="https://img.shields.io/badge/Hugging%20Face-FFusion--BaSE-blue" alt="Hugging Face Model"></a> |
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<a href="https://github.com/1e-2" target="_new" rel="ugc"><img src="https://img.shields.io/badge/GitHub-1e--2-green" alt="GitHub"></a> |
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<a href="https://www.facebook.com/FFusionAI/" target="_new" rel="ugc"><img src="https://img.shields.io/badge/Facebook-FFusionAI-blue" alt="Facebook"></a> |
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<a href="https://civitai.com/models/82039/ffusion-ai-sd-21" target="_new" rel="ugc"><img src="https://img.shields.io/badge/Civitai-FFusionAI-blue" alt="Civitai"></a> |
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</div> |
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### 🧨 Diffusers |
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Make sure to upgrade diffusers to >= 0.18.0: |
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``` |
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pip install diffusers --upgrade |
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``` |
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In addition make sure to install `transformers`, `safetensors`, `accelerate` as well as the invisible watermark: |
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``` |
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pip install invisible_watermark transformers accelerate safetensors |
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``` |
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You can use the model then as follows |
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```py |
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from diffusers import DiffusionPipeline |
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import torch |
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pipe = DiffusionPipeline.from_pretrained("FFusion/FFusionXL-09-SDXL", torch_dtype=torch.float16, use_safetensors=True, variant="fp16") |
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pipe.to("cuda") |
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# if using torch < 2.0 |
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# pipe.enable_xformers_memory_efficient_attention() |
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prompt = "An astronaut riding a green horse" |
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images = pipe(prompt=prompt).images[0] |
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``` |
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When using `torch >= 2.0`, you can improve the inference speed by 20-30% with torch.compile. Simple wrap the unet with torch compile before running the pipeline: |
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```py |
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pipe.unet = torch.compile(pipe.unet, mode="reduce-overhead", fullgraph=True) |
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``` |
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If you are limited by GPU VRAM, you can enable *cpu offloading* by calling `pipe.enable_model_cpu_offload` |
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instead of `.to("cuda")`: |
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```diff |
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- pipe.to("cuda") |
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+ pipe.enable_model_cpu_offload() |
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``` |
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## Uses |
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### Direct Use |
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The model is intended for research purposes only. Possible research areas and tasks include |
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- Generation of artworks and use in design and other artistic processes. |
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- Applications in educational or creative tools. |
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- Research on generative models. |
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- Safe deployment of models which have the potential to generate harmful content. |
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- Probing and understanding the limitations and biases of generative models. |
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Excluded uses are described below. |
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### Out-of-Scope Use |
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The model was not trained to be factual or true representations of people or events, and therefore using the model to generate such content is out-of-scope for the abilities of this model. |
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## Limitations and Bias |
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### Limitations |
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- The model does not achieve perfect photorealism |
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- The model cannot render legible text |
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- The model struggles with more difficult tasks which involve compositionality, such as rendering an image corresponding to “A red cube on top of a blue sphere” |
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- Faces and people in general may not be generated properly. |
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- The autoencoding part of the model is lossy. |
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### Bias |
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While the capabilities of image generation models are impressive, they can also reinforce or exacerbate social biases. |
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**Attribution:** |
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"SDXL 0.9 is licensed under the SDXL Research License, Copyright (c) Stability AI Ltd. All Rights Reserved." |
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## License |
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[SDXL 0.9 Research License](https://huggingface.co/stabilityai/stable-diffusion-xl-base-0.9/blob/main/LICENSE.md)" |
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[FFXL 0.9 Research License](https://huggingface.co/FFusion/FFusionXL-09-SDXL/blob/main/LICENSE.md)" |
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[](mailto:[email protected]) |
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## SAMPLES |
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