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https://www.keyword-suggest-tool.com/search/vertex+to+standard+quadratic+form+calculator/
[ "# Vertex to standard quadratic form calculator\n\nVertex to standard quadratic form calculator keyword after analyzing the system lists the list of keywords related and the list of websites with related content, in addition you can see which keywords most interested customers on the this website\n\n## Keyword Suggestions\n\nStandard to vertex quadratic form calculator\n\n( Please select at least 2 keywords )\n\n#### Most Searched Keywords\n\n## Websites Listing\n\nWe found at least 10 Websites Listing below when search with vertex to standard quadratic form calculator on Search Engine\n\n### Quadratic Equation Calculator - Symbolab\n\nDA: 16 PA: 37 MOZ Rank: 53\n\n### Functions Vertex Calculator - Symbolab\n\nFree functions vertex calculator - find function's vertex step-by-step. This website uses cookies to ensure you get the best experience. ... Arithmetic Mean Geometric Mean Quadratic Mean Median Mode Order Minimum Maximum Probability Mid-Range Range Standard Deviation Variance Lower Quartile Upper Quartile ... just click the link in the email we ...\n\nhttps://www.symbolab.com/solver/function-vertex-calculator\n\nDA: 16 PA: 34 MOZ Rank: 50\n\nOur quadratic equation formula solver is designed to solve all types of quadratic equations. Be it standard form, factored form or vertex form, you will get your answer within a few minutes. Works on all browsers. The best part of the online quadratic equation calculator is it is compatible with all type of window versions as well as browsers.\n\nDA: 20 PA: 31 MOZ Rank: 51\n\n### Vertex form calculator - Cannabisser\n\nWhy not try this out? Thousands of users are using our software to conquer their algebra homework. Instead of x², you can also write x^2. Contact - They are pretty smart and first class . Now that's customer service! From vertex form to standard form. However, I asked my son to give it a try. If the coefficient of the term x2 is positive, the vertex will be in the bottom of the U- shaped curve.\n\nhttps://www.cannabisser.com/site/0aff5d-vertex-form-calculator\n\nDA: 19 PA: 35 MOZ Rank: 54\n\n### Quadratic equations vertex and standard form online ...\n\nSte. 105-181 19179 Blanco Rd #181 San Antonio, TX 78258 USA\n\nDA: 12 PA: 50 MOZ Rank: 91\n\n### Braingenie | Changing from Vertex to Standard Form\n\nIdentifying the Vertex as Max or Min Given Intercept Form Finding the Vertex Given Intercept Form Graphing y = a(x + p)(x + q) Using the Vertex and X-Intercepts Solving Word Problems Using Intercept Form Changing Between Vertex, Intercept and Standard Form\n\nhttps://braingenie.ck12.org/skills/106803\n\nDA: 19 PA: 14 MOZ Rank: 33\n\n### Relations and Functions Calculators\n\nEvaluate Functions. This calculator will: (1.) Evaluate a function for a specified value. (2.) Return the answer in the simplest form. (3.) Graph the function and indicate the specified value.\n\nhttps://relations-functions.appspot.com/calculators.html\n\nDA: 31 PA: 17 MOZ Rank: 48\n\n### Solver Convert to Vertex Form and Graph\n\nThis Solver (Convert to Vertex Form and Graph) was created by by ccs2011(207) : View Source, Show, Put on YOUR site About ccs2011: Convert to Vertex Form and Graph. Enter quadratic equation in standard form:--> x 2 + x + This solver has been accessed 2388320 times. ...\n\nDA: 15 PA: 50 MOZ Rank: 68\n\n### Mod 2 notes.pdf - MAT 110 Module 2 Notes Module 2 Notes ...\n\nMAT 110 - Module 2 Notes Standard to Vertex Form 1. When a quadratic function is in the form y = ax 2 + bx + c, to change is to vertex form: (a) Calculate the axis of symmetry using the formula x =-b 2 a (b) When you find this value, plug it back into the function. (c) the a value does not change (d) Write the equation in vertex for y = a (x-h ...\n\nhttps://www.coursehero.com/file/77791204/mod-2-notespdf/\n\nDA: 18 PA: 30 MOZ Rank: 48\n\n### Solver Completing the Square to Get a Quadratic into ...\n\nThis Solver (Completing the Square to Get a Quadratic into Vertex Form) was created by by jim_thompson5910(35256) : View Source, Show, Put on YOUR site About jim_thompson5910: If you need more math help, then you can email me. I charge \\$2 for steps, or \\$1 for answers only.\n\nhttps://www.algebra.com/algebra/homework/Polynomials-and-rational-expressions/completing-the-square.solver\n\nDA: 15 PA: 50 MOZ Rank: 98" ]
[ null ]
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https://rdrr.io/cran/tipmap/f/inst/doc/introduction.Rmd
[ "# Introduction to the R package 'tipmap'\" In tipmap: Tipping Point Analysis for Bayesian Dynamic Borrowing\n\n```knitr::opts_chunk\\$set(\necho = T, collapse = T, warning = F, message = F,\nprompt = T, comment = \"#\",\nout.width = \"100%\"\n)\n```\n\n## Purpose of the package\n\nThe R package tipmap implements a tipping point analysis for clinical trials that employ Bayesian dynamic borrowing of a treatment effect from external evidence via robust meta-analytic predictive (MAP) priors. A tipping point analysis allows to assess how much weight on the informative component of a robust MAP prior is needed to conclude that the investigated treatment is efficacious, based on the total evidence. The package mainly provides an implementation of a graphical approach proposed by (Best et al., 2021) for different one-sided evidence levels (80%, 90%, 95%, 97.5%).\n\nTipping point analyses can be useful both at the planning and the analysis stage of a trial that uses external information. At the planning stage, they can help determining (pre-specifying) a weight of the informative component of the MAP prior for the primary analysis. Various possible results of the planned trial in the target population and the consequences for the treatment effect estimate and statistical inferences based on the total evidence may be explored for a range of weights. Through this exercise, in addition to other criteria, decision-makers can develop a sense under which circumstances they would still feel comfortable to establish efficacy with a specific level of certainty. A preferred weight can then rather straightforwardly be derived. The tipping point analysis can thus help to elicit sensible weights. At the analysis stage, tipping point analyses can be used as sensitivity analysis to assess the dependency of the treatment effect estimate and statistical inferences on weight of the informative component of the MAP prior. This may also be understood in the sense of a reverse-Bayes analysis (Held et al., 2022).\n\nThis vignette shows an exemplary application of the tipping point analysis with hypothetical data.\n\nFurther functions in this package (not illustrated in this vignette) facilitate the specification of a robust MAP prior via expert elicitation, specifically the choice of a primary weight, using the roulette method, as well as the computation of the posterior distribution of the treatment effect with stochastic weights elicited from a panel of experts.\n\nIntended use of the tipmap-package is the planning, analysis and interpretation of (small) clinical trials in pediatric drug development, where extrapolation of efficacy, often through Bayesian methods, has become increasingly common (Gamalo et al., 2022; ICH, 2022). However, the applicability of the package is generally wider.\n\nFor the implementation of the MAP prior approach, including computation of the posterior distribution, the RBesT-package is used (Weber et al., 2021).\n\n## Data generation for an exemplary tipping point analysis\n\nIn this vignette, we assume that results from three clinical trials conducted in adult patients (the source population) are available, which share key features with a new trial among pediatric patients (the target population). For example, they had been conducted in the same indication, studied the same drug and provided results on an endpoint of interest for the target population. This means a certain degree of exchangeability between the trials in the source and target population can be assumed. The similarity in disease and response to treatment between source and target population always need to be carefully considered in any setting, usually by clinical experts in the disease area.\n\nWe assume that it is supported by medical evidence and now planned to consider the trials in adult patients in a Bayesian dynamic borrowing approach, and we would like to create a robust MAP prior (Schmidli et al., 2014). The treatment effect measure of interest is assumed to be a mean difference between a treated group and a control group with respect to a continuous endpoint.\n\n### Derivation of MAP prior based on trials in the source population\n\nWe start by specifying an object that contains the prior data.\n\nThe function `create_prior_data()` takes vectors of total sample sizes, treatment effect estimates and their standard errors as arguments and generates a data frame. A study label is optional.\n\n```library(tipmap)\nprior_data <- create_prior_data(\nn_total = c(160, 240, 320),\nest = c(1.16, 1.43, 1.59),\nse = c(0.46, 0.35, 0.28)\n)\n```\n```# compute standard deviation of change in each arm (assumed equal);\n# for two-sample data:\nsd <- prior_data\\$se / sqrt(1/(prior_data\\$n_total/2) + 1/(prior_data\\$n_total/2))\nsigma1 <- mean(sd)\nsigma1 # = 2.70826\nnt <- 15; nc <- 15\nf <- (nt+nc)^2 / (nt*nc)\nsigma2 <- sqrt(f)*sigma1\nsigma2 # = 5.41652\n```\n```print(prior_data)\n```\n\nWe then generate a MAP prior from our prior data using the RBesT-package (Weber et al., 2021).\n\n```set.seed(123)\nuisd <- 5.42\nmap_mcmc <- RBesT::gMAP(\nformula = cbind(est, se) ~ 1 | study_label,\ndata = prior_data,\nfamily = gaussian,\nweights = n_total,\ntau.dist = \"HalfNormal\",\ntau.prior = cbind(0, uisd / 16),\nbeta.prior = cbind(0, uisd)\n)\n```\n\nA few additional specifications are needed to be made to fit the MAP prior model; for details see Neuenschwander and Schmidli (2020) or Weber et al. (2021). The variable `uisd` here represents an assumed unit-information standard deviation and the specification of the prior on between-trial heterogeneity parameter tau follows recommendations to consider moderate heterogeneity for a two-group parameter, such as the mean difference (Neuenschwander and Schmidli, 2020).\n\nThis is a summary of the fitted model based on samples from the posterior distribution:\n\n```summary(map_mcmc)\n```\n\nA forest plot of the Bayesian meta-analysis is shown in Figure 1. It is augmented with meta-analytic shrinkage estimates per trial. The figure shows the per-trial point estimates (light point) and the 95% frequentist confidence intervals (dashed line) and the model derived median (dark point) and the 95% credible interval of the meta-analytic model.\n\n```plot(map_mcmc)\\$forest_model\n```\n\nSubsequently, the MAP prior is approximated by a mixture of conjugate normal distributions. The parametric form facilitates the computation of posteriors when the MAP prior is combined with results from the trial in the target population.\n\n```map_prior <- RBesT::automixfit(\nsample = map_mcmc,\nNc = seq(1, 4),\nk = 6,\nthresh = -Inf\n)\n```\n\nThe approximation yields a mixture of two normals:\n\n```print(map_prior)\n```\n\nThe density of the parametric mixture together with a histogram of MCMC samples from the `map_mcmc` object is shown in Figure 2.\n\n```plot(map_prior)\\$mix\n```\n\nThe derivation of the MAP prior is now complete. For normal likelihoods the parametric representation by a mixture of normals can be used to calculate posterior distributions analytically.\n\n### Trial results in the target population\n\nWe now create a numeric vector with data on pediatric trial (the total sample size, the treatment effect estimate and its standard error). In the planning stage, this may be an expected result.\n\n```pediatric_trial <- create_new_trial_data(n_total = 30, est = 1.02, se = 1.4)\n```\n```print(pediatric_trial)\n```\n\nThe function `create_new_trial_data()` computes quantiles, assuming normally distributed errors. This is merely used to plot a confidence interval for the treatment effect estimate obtained in the target trial in the tipping point plot.\n\n## Performing the tipping point analysis\n\n### Computation of posteriors for a range of weights\n\nWe can now compute posterior distributions for a range of weights using the function `create_posterior_data()`.\n\n```posterior <- create_posterior_data(\nmap_prior = map_prior,\nnew_trial_data = pediatric_trial,\nsigma = uisd)\n```\n```head(posterior, 4)\n```\n```tail(posterior, 4)\n```\n\nThe resulting data frame has `r dim(posterior)` rows and `r dim(posterior)` columns. The weights increase incrementally in steps of 0.005 from 0 to 1, i.e.\\ posterior quantiles for `r length(posterior\\$weight)` weights are computed. For each weight the data frame contains the following `r dim(posterior)-1` posterior quantiles.\n\n```length(posterior\\$weight)\ndim(posterior)\nclass(posterior)\ncolnames(posterior)[-1]\n```\n```colnames(posterior)[-1]\n```\n\nThese posterior quantiles can be directly used for inferences based on the total evidence (new data and prior combined). They reflect one-sided 99%, 97.5%, 95%, 90%, 80%, and 50% evidence levels for a given weight, respectively.\n\n### Creating the tipping point plot\n\nThe function to produce the tipping point plot is called `tipmap_plot()`, it requires a dataframe with data on all components generated by the function `create_tipmap_data()`.\n\n```tipmap_data <- create_tipmap_data(\nnew_trial_data = pediatric_trial,\nposterior = posterior,\nmap_prior = map_prior)\n```\n```(p1 <- tipmap_plot(tipmap_data = tipmap_data))\n```\n\nIn the center of the plot, a funnel-shaped display of quantiles of the posterior distribution (reflecting one-sided evidence-levels) is shown for given weights of the informative component of the MAP prior. The intersections between the lines connecting the respective quantiles and the horizontal line at 0 (the null effect) are referred to as tipping points (indicated by vertical lines in red color). They indictae the minimum weight that is required to conclude that the treatment is efficacious for a given one-sided evidence level (Best et al., 2021). On the left and right side of the plot, the treatment effect estimate obtained in the trial in the (pediatric) target population (with 95% confidence interval) and the MAP prior (with 95% credible interval) are shown, respectively.\n\nThe plot is a `ggplot`-object that can be modified accordingly. For example, if we had chosen a primary weight of 0.38, we could add a vertical reference line at this position. There are additional features to customize the plot in the `tipmap_plot()` function, see `help(tipmap_plot)`.\n\n```primary_weight <- 0.38\n(p2 <- p1 + ggplot2::geom_vline(xintercept = primary_weight, col=\"green4\"))\n```\n\nWe see from Figure 4 that, for a weight of 0.38, there is a probability of larger than 90% but less than 95% based on the posterior distribution that the treatment effect is larger than 0, i.e. the treatment is efficacious.\n\n### Extracting quantities of interest\n\nThe data frame with posteriors for all weights can be filtered to obtain posterior quantiles for weights of specific interest by the function `get_posterior_by_weight()`:\n\n```get_posterior_by_weight(\nposterior = posterior,\nweight = c(primary_weight)\n)\n```\n\nThe function `get_tipping_points()` extracts tipping points for one-sided 80%, 90%, 95% and 97.5% evidence levels, respectively.\n\n```tipp_points <- get_tipping_points(\ntipmap_data = tipmap_data,\nquantile = c(0.2, 0.1, 0.05, 0.025)\n)\ntipp_points\n```\n\nCalculating the precise posterior probability that treatment effect exceeds a threshold value is possible via functions in the RBesT-package.\n\n```prior_primary <- RBesT::robustify(\npriormix = map_prior,\nweight = (1 - primary_weight),\nm = 0,\nn = 1,\nsigma = uisd\n)\n```\n```posterior_primary <- RBesT::postmix(\npriormix = prior_primary,\nm = pediatric_trial[\"mean\"],\nse = pediatric_trial[\"se\"]\n)\n```\n```summary(posterior_primary)\n```\n\nThe posterior probability that the treatment effect is larger than 0, 0.5 and 1, respectively, can be assessed through the cumulative distribution function of the posterior.\n\n```round(1 - RBesT::pmix(posterior_primary, q = 0), 3)\nround(1 - RBesT::pmix(posterior_primary, q = 0.5), 3)\nround(1 - RBesT::pmix(posterior_primary, q = 1), 3)\n```\n\nThis is illustrated by a cumulative density curve of the posterior.\n\n```library(ggplot2)\nplot(posterior_primary, fun = RBesT::pmix) +\nscale_x_continuous(breaks = seq(-1, 2, 0.5)) +\nscale_y_continuous(breaks = 1-c(1, 0.927, 0.879, 0.782, 0.5, 0),\nlimits = c(0,1),\nexpand = c(0,0)\n) +\nylab(\"Cumulative density of posterior with w=0.38\") +\nxlab(\"Quantile\") +\ngeom_segment(aes(x = 0,\ny = RBesT::pmix(mix = posterior_primary, q = 0),\nxend = 0,\nyend = 1),\ncol=\"red\") +\ngeom_segment(aes(x = 0.5,\ny = RBesT::pmix(mix = posterior_primary, q = 0.5),\nxend = 0.5,\nyend = 1),\ncol=\"red\") +\ngeom_segment(aes(x = 1,\ny = RBesT::pmix(mix = posterior_primary, q = 1),\nxend = 1,\nyend = 1),\ncol=\"red\") +\ntheme_bw()\n```\n\nAs a further example, for the weight corresponding to the tipping point of the one-sided evidence-level of 95% (=0.51), we would obtain a posterior probability of 95% that the treatment effect is larger than 0.\n\n```tipp_points\n```\n```prior_95p <- RBesT::robustify(\npriormix = map_prior,\nweight = (1 - tipp_points),\nm = 0,\nn = 1,\nsigma = uisd\n)\n```\n```posterior_95p <- RBesT::postmix(\npriormix = prior_95p,\nm = pediatric_trial[\"mean\"],\nse = pediatric_trial[\"se\"]\n)\n```\n```round(1 - RBesT::pmix(posterior_95p, q = 0), 3)\n```\n\n## References\n\nBest, N., Price, R. G., Pouliquen, I. J., et al. (2021) Assessing efficacy in important subgroups in confirmatory trials: An example using Bayesian dynamic borrowing. Pharm Stat, 20, 551–562. DOI: 10.1002/pst.2093\n\nGamalo, M., Bucci-Rechtweg, C., Nelson, R. M., et al. (2022) Extrapolation as a Default Strategy in Pediatric Drug Development. Ther Innov Regul Sci, 56, 883–894. DOI: 10.1007/s43441-021-00367-9\n\nHeld, L., Matthews, R., Ott, M., et al. (2022) Reverse-Bayes methods for evidence assessment and research synthesis. Res Synth Methods, 13, 295–314. DOI: 10.1002/jrsm.1538\n\nICH (2022) International Council for Harmonisation of Technical Requirements for Pharmaceuticals for Human Use (ICH) guideline E11A on pediatric extrapolation. EMA/CHMP/ICH/205218/2022.\n\nNeuenschwander, B. and Schmidli, H. (2020) Use of Historical Data. In Bayesian Methods in Pharmaceutical Research (eds E. Lesaffre, G. Baio, and B. Boulanger), pp. 111–129. 1st ed. Chapman & Hall/ CRC Biostatistics Series. DOI: 10.1201/9781315180212\n\nSchmidli, H., Gsteiger, S., Roychoudhury, S., et al. (2014) Robust meta-analytic-predictive priors in clinical trials with historical control information. Biometrics, 70, 1023–1032. DOI: 10.1111/biom.12242\n\nWeber, S., Li, Y., Seaman III, J. W., et al. (2021) Applying Meta-Analytic-Predictive Priors with the R Bayesian Evidence Synthesis Tools. J Stat Softw, 100, 1–32. DOI: 10.18637/jss.v100.i19\n\n## Try the tipmap package in your browser\n\nAny scripts or data that you put into this service are public.\n\ntipmap documentation built on Dec. 8, 2022, 1:13 a.m." ]
[ null ]
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https://biology.reachingfordreams.com/chemistry-cheat-sheet/chemical-equilibrium/40-changes-in-equilibrium-systems
[ "Le Châtelier’s principle - a generalization that states that:\n\nchemical systems at equilibrium shift to restore equilibrium when a change occurs that disturbs the equilibrium.\n\nAn adjustment by a system at equilibrium that results in a change in the concentrations of reactants and products is called an equilibrium shift.\n\nLe Châtelier’s principle allows chemists to predict the qualitative effects of changes in concentration, pressure, and temperature on a chemical reaction system at equilibrium.\n\nWe can tell a reaction is at equilibrium if the reaction quotient (Q) is equal to the equilibrium constant (K).\n\nIf a system at equilibrium is subjected to a perturbance or stress (such as a change in concentration) the position of equilibrium changes.\n\nSince this stress affects the concentrations of the reactants and the products, the value of Q will no longer equal the value of K.\n\nTo re-establish equilibrium, the system will either shift toward the products (if Q < K) or the reactants (if Q > K) until Q returns to the same value as K.\n\n## The Relationship Between the Equilibrium Constant and the Reaction Quotient", null, "If reaction quotient (Qc) is greater than equilibrium constant (Kc), the numerator of the reaction quotient expression must be very large.\n\nThe concentrations of the chemicals on the right side of the equation\n\n(mA + nB + ⇌ xC + yD)\n\nmust be more than their concentrations at equilibrium.\n\nIn this situation, the system attains equilibrium by moving to the left.\n\nConversely, if Qc is less than Kc , the system attains equilibrium by moving to the right.\n\n## Effect of Change in Concentration on Equilibrium", null, "(a) The test tube contains 0.1 M Fe3+.\n\n(b) Thiocyanate ion has been added to solution in (a), forming the red Fe(SCN)2+ ion.\n\nFe3+ (aq) + SCN (aq) ⇌ Fe(SCN)2+ (aq).\n\n(c) Silver nitrate has been added to the solution in (b), precipitating some of the SCN − as the white solid AgSCN.\n\nAg+ (aq) + SCN (aq) ⇌ AgSCN(s).\n\nThe decrease in the SCN concentration shifts the first equilibrium in the solution to the left, decreasing the concentration (and lightening color) of the Fe(SCN)2+.\n\nThe stress on the system in the image above is the reduction of the equilibrium concentration of SCN (lowering the concentration of one of the reactants would cause Q to be larger than K).\n\nAs a consequence, Le Châtelier's principle leads us to predict that the concentration of Fe(SCN)2+ should decrease, increasing the concentration of SCN part way back to its original concentration, and increasing the concentration of Fe3+ above its initial equilibrium concentration.\n\n### Collision Theory and Concentration Changes in an Equilibrium System\n\nAccording to collision theory, entities in a chemical system must collide to react.\n\nWhen the concentration of an entity in a chemical reaction system is increased, it is more likely that that entity will collide with other entities. There are simply more of them present. However, only collisions between reactant entities can potentially contribute to a chemical reaction. Even so, the more frequently collisions occur overall, the more likely it is that a chemical reaction will take place.\n\nCollision theory explains the response of a chemical reaction system at equilibrium to a change in concentration as the result of random collisions and probability.\n\nWhen we add more reactant entities in an equilibrium system, the equilibrium shifts to the right because the number of successful collisions for the forward reaction increases.\n\nIf, instead, we add more product entities, then the number of successive collisions for the reverse reaction will increase and the equilibrium will shift to the left.\n\n## Le Châtelier’s Principle and Changes in Energy\n\nFor a given system at equilibrium, the value of the equilibrium constant depends only on temperature.\n\nChanging the temperature of a reacting mixture changes the rate of the forward and reverse reactions by different amounts, because the forward and reverse reactions have different activation energies.\n\nA reacting mixture at one temperature has an equilibrium constant whose value changes if the mixture is allowed to reach equilibrium at a different temperature.\n\nTo apply Le Châtelier’s principle and predict how a change in energy will affect a chemical system at equilibrium, we can think of energy as a reactant or a product.\n\nEnergy is absorbed in an endothermic reaction. If we consider energy to be a reactant, we can write the word equation:\n\nreactants + energy ⇌ products\n\nSince energy is released during an exothermic reaction, we can consider energy to be a reactant and write\n\nreactants ⇌ products + energy\n\nTo use Le Châtelier’s principle to predict how an equilibrium will shift in response to a change in energy, consider how the system can counteract this shift.\n\n### Endothermic Reactions\n\nIf an endothermic reaction is cooled (thermal energy removed), we can consider that the quantity of one of the reactants has been decreased.\n\nWe can therefore predict that the equilibrium will shift to the left (toward the reactants), and energy will be released.\n\nIf thermal energy were added to this equilibrium system by heating, the equilibrium would likely shift to the right.\n\n### Exothermic Reactions\n\nIf thermal energy is removed from an exothermic reaction—where energy is a product—then the equilibrium will shift to the right (toward the products), and energy will be released to counteract the change.\n\nIf energy is added to an exothermic reaction, the equilibrium will shift to the left to compensate for the change, and the energy will be used as products are converted to reactants.\n\n## Le Châtelier’s Principle and Changes in Gas Volume\n\nSometimes we can change the position of equilibrium by changing the pressure of a system. However, changes in pressure have a measurable effect only in systems in which gases are involved, and then only when the chemical reaction produces a change in the total number of gas molecules in the system.\n\nAn easy way to recognize such a system is to look for different numbers of moles of gas on the reactant and product sides of the equilibrium.\n\nWhile evaluating pressure (as well as related factors like volume), it is important to remember that equilibrium constants are defined with regard to concentration (for Kc) or partial pressure (for Kp).\n\nSome changes to total pressure, like adding an inert gas that is not part of the equilibrium, will change the total pressure but not the partial pressures of the gases in the equilibrium constant expression.\n\nThus, addition of a gas not involved in the equilibrium will not perturb the equilibrium.\n\nAs we increase the pressure of a gaseous system at equilibrium, either by decreasing the volume of the system or by adding more of one of the components of the equilibrium mixture, we introduce a stress by increasing the partial pressures of one or more of the components.\n\nIn accordance with Le Châtelier's principle, a shift in the equilibrium that reduces the total number of molecules per unit of volume will be favored because this relieves the stress.\n\nThe reverse reaction would be favored by a decrease in pressure.\n\nConsider what happens when we increase the pressure on a system in which N2, H2, and NH3 are at equilibrium:\n\nN(g) + 3H2 (g) ⇌ 2NH3 (g)\n\nThe formation of additional amounts of NH3 decreases the total number of molecules in the system because each time two molecules of NH3 form, a total of four molecules of N2 and H2 are consumed.\n\nThis reduces the total pressure exerted by the system and reduces, but does not completely relieve, the stress of the increased pressure.\n\nOn the other hand, a decrease in the pressure on the system favors decomposition of NH3 into N2 and H2 , which tends to restore the pressure.\n\nLe Châtelier’s Principle and Changes in Gas Volume", null, "(a) The container holds the equilibrium reaction N(g) + 3H2 (g) ⇌ 2NH3 (g) .\n\n(b) The volume of the container is decreased by half, which doubles the total pressure of the system.\n\n(c) The reaction shifts to the right, which decreases the total number of particles in the container and the total pressure by producing more ammonia, NH3 (g).\n\n## The Effect of a Catalyst on Equilibrium\n\nA catalyst speeds up the rate of a reaction, either by allowing a different reaction mechanism or by providing additional mechanisms.\n\nThe overall effect is to lower the activation energy, which increases the rate of reaction.\n\nThe activation energy is lowered the same amount for the forward and reverse reactions, however. There is the same increase in reaction rates for both reactions.\n\nAs a result, a catalyst does not affect the position of equilibrium. It only affects the time that is taken to achieve equilibrium.\n\nThe Effects of Changing Conditions on a System at Equilibrium\n\n Type of reaction Change to system Effect on Kc Direction of change all reactions increasing any reactant concentration, or decreasing any product concentration no effect toward products decreasing any reactant concentration, or increasing any product concentration no effect toward reactants using a catalyst no effect no change exothermic increasing temperature decreases toward reactants decreasing temperature increases toward products endothermic increasing temperature increases toward products decreasing temperature decreases toward reactants equal number of reactant and product gas molecules changing the volume of the container, or adding a non-reacting gas no effect no change more gaseous product molecules than reactant gaseous molecules decreasing the volume of the container at constant temperature no effect toward reactants increasing the volume of the container at constant temperature, or adding a non-reacting gas at contstant pressure no effect toward products fewer gaseous product molecules than reactant gaseous molecules decreasing the volume of the container at constant temperature no effect toward products increasing the volume of the container at constant temperature no effect toward reactants" ]
[ null, "https://biology.reachingfordreams.com/plugins/system/lazyloadforjoomla/src/assets/images/blank.gif", null, "https://biology.reachingfordreams.com/plugins/system/lazyloadforjoomla/src/assets/images/blank.gif", null, "https://biology.reachingfordreams.com/plugins/system/lazyloadforjoomla/src/assets/images/blank.gif", null ]
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https://ux.stackexchange.com/questions/140997/when-to-use-white-text-on-top-of-a-color-for-readability/140998
[ "# When to use white text on top of a color for readability?\n\nI am setting the background of a textbox to some variable color.\n\nI know the hex/rgb of the color I'm setting the background to.\n\nNormally the text is black (or very very dark gray).\n\nOn some colors, it would be far easier to read the text if it were white instead of black.\n\nHow can it be programmatically decided, given a hex/rgb color, whether white or black text would be more readable if placed on top of it?\n\nThis is standardised by W3C's Web Content Accessibility Guidelines (WCAG) and decided based on the contrast ratio between colors, which you can calculate programmatically.\n\nIn your case, you would want to use whichever of white or black has a higher contrast ratio with the background color.\n\nHere is just one of many articles on how to calculate contrast ratio in JavaScript: dev.to: Building your own color contrast checker.\n\nFrom that article:\n\nThe contrast ratio between two colors is calculated using a mathematical function that we will see below. The resulting value of that calculation will go from 1 to 21 (normally represented as 1:1 or 21:1 respectively). If both are the exact same color, the contrast ratio will be 1; The highest contrast is between white and black, and exact opposites in the color wheel (complementary colors) will have high values too.\n\nThe mathematical formula is documented by WCAG.\n\ncontrast ratio\n(L1 + 0.05) / (L2 + 0.05), where\n\n• L1 is the relative luminance of the lighter of the colors, and\n• L2 is the relative luminance of the darker of the colors.\n\nIt uses relative luminance:\n\nFor the sRGB colorspace, the relative luminance of a color is defined as L = 0.2126 * R + 0.7152 * G + 0.0722 * B where R, G and B are defined as:\n\n• if RsRGB <= 0.03928 then R = RsRGB/12.92 else R = ((RsRGB+0.055)/1.055) ^ 2.4\n• if GsRGB <= 0.03928 then G = GsRGB/12.92 else G = ((GsRGB+0.055)/1.055) ^ 2.4\n• if BsRGB <= 0.03928 then B = BsRGB/12.92 else B = ((BsRGB+0.055)/1.055) ^ 2.4\n\nand RsRGB, GsRGB, and BsRGB are defined as:\n\n• RsRGB = R8bit/255\n• GsRGB = G8bit/255\n• BsRGB = B8bit/255\n\nThere exist free online contrast calculators like this one from WebAIM. In your case, add #fff as one value and your other color as the other. In most cases, you only need WCAG AA passing to be considered accessible. WCAG AA considers a 4.5:1 ratio to be passing. Usually it's good to go just a little darker than a minimum passing value, and the contrast checker can help you find a value for that.\n\n• 4.5:1 is passing, but 7:1 is the recommendation for main body text (for example, the text of the answers here on SE), and speaking from experience as someone who actually has vision issues affected by this, ‘passing’ should not be considered ‘good enough’ unless you are dealing with large font sizes or exceptionally good fonts. Sep 14, 2021 at 11:52" ]
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https://se.mathworks.com/help/control/ref/dlyapchol.html
[ "# dlyapchol\n\nSquare-root solver for discrete-time Lyapunov equations\n\n## Syntax\n\n```R = dlyapchol(A,B) X = dlyapchol(A,B,E) ```\n\n## Description\n\n`R = dlyapchol(A,B)` computes a Cholesky factorization `X = R'*R` of the solution `X` to the Lyapunov matrix equation:\n\n```A*X*A'- X + B*B' = 0 ```\n\nAll eigenvalues of `A` matrix must lie in the open unit disk for `R` to exist.\n\n`X = dlyapchol(A,B,E)` computes a Cholesky factorization `X = R'*R` of `X` solving the Sylvester equation\n\n```A*X*A' - E*X*E' + B*B' = 0 ```\n\nAll generalized eigenvalues of (`A`,`E`) must lie in the open unit disk for `R` to exist.\n\n## Algorithms\n\n`dlyapchol` uses SLICOT routines SB03OD and SG03BD.\n\n## References\n\n Bartels, R.H. and G.W. Stewart, \"Solution of the Matrix Equation AX + XB = C,\" Comm. of the ACM, Vol. 15, No. 9, 1972.\n\n Hammarling, S.J., “Numerical solution of the stable, non-negative definite Lyapunov equation,” IMA J. Num. Anal., Vol. 2, pp. 303-325, 1982.\n\n Penzl, T., ”Numerical solution of generalized Lyapunov equations,” Advances in Comp. Math., Vol. 8, pp. 33-48, 1998." ]
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https://forum.inductiveautomation.com/t/best-way-to-construct-timestamp-from-elements/11919
[ "# Best way to construct timestamp from elements\n\nI want to construct a timestamp from individual elements (year, month, day, hour, minute, second).\nThe elements are initially coming from a PLC and then I’m writing into a database.\n\nThis script (admittedly with some hard-coded values) appears to work:\n\n[code]from datetime import datetime\ndt_y = 2016\ndt_m = 1\ndt_d = 20\ndt_h = 12\ndt_mn = 34\ndt_s = 56\ntstamp=datetime(dt_y, dt_m, dt_d, dt_h, dt_mn, dt_s)\n\nsql = “UPDATE mytable SET t_stamp=? WHERE ID=1”\nsystem.db.runPrepUpdate(sql,[tstamp])[/code]\n\nBut is this the best way of doing it?\n\nThere’s probably not a “right” answer to this. What you have is certainly acceptable. Since Ignition is built with java and jython has to convert to java objects to move data through the system, I tend to use the corresponding java objects directly. There are various differences, but the same end effect. Somewhat more efficient if needed in big loops.[code]from java.util import Date\ndt_y = 2016\ndt_m = 1\ndt_d = 20\ndt_h = 12\ndt_mn = 34\ndt_s = 56\ntstamp=Date(dt_y-1900, dt_m-1, dt_d, dt_h, dt_mn, dt_s)\n\nsql = “UPDATE mytable SET t_stamp=? WHERE ID=1”\nsystem.db.runPrepUpdate(sql,[tstamp])[/code]" ]
[ null ]
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https://lapsedgeographer.london/2020-04/case_when/
[ "# When's the case for case_when() ?", null, "Last week I was taking some colleagues through the code for my COVID19 PDF scraping and afterwards one sent me a message asking about a chunk of code that used the `dplyr::case_when()` function. In particular they wanted to know why `case_when()` uses the tilde (`~`)1, which led to a bit more of a generalised conversation about `case_when()` and how it works.\n\nIn your script, when you assign the values to the entity and position columns inside your `case_when()` call why do you use `~` rather than `=` or `<-`? It felt too specific and potentially obvious to others to ask during your talk. I understand why it’s used in model building for regression etc… but not really in this context.\n\nThis is a great question and formulating the answer actually helped crystallise something I knew about R but hadn’t really processed before then.\n\n``````subnational_datapoints <- subnational_data %>%\nfilter(y == 36 | y == 104 | y == 242 |\ny == 363 | y == 431 | y == 568) %>%\nmutate(\nentity = case_when(\ny == 36 ~ \"location\",\ny == 104 & x == 36 ~ \"retail_recr\",\ny == 104 & x == 210 ~ \"grocery_pharm\",\ny == 104 & x == 384 ~ \"parks\",\ny == 242 & x == 36 ~ \"transit\",\ny == 242 & x == 210 ~ \"workplace\",\ny == 242 & x == 384 ~ \"residential\",\ny == 363 ~ \"location\",\ny == 431 & x == 36 ~ \"retail_recr\",\ny == 431 & x == 210 ~ \"grocery_pharm\",\ny == 431 & x == 384 ~ \"parks\",\ny == 568 & x == 36 ~ \"transit\",\ny == 568 & x == 210 ~ \"workplace\",\ny == 568 & x == 384 ~ \"residential\"),\nposition = case_when(\ny == 36 ~ \"first\",\ny == 104 ~ \"first\",\ny == 242 ~ \"first\",\ny == 363 ~ \"second\",\ny == 431 ~ \"second\",\ny == 568 ~ \"second\")\n)\n``````\n\n## The basics\n\nThe documentation for `dplyr::case_when()` calls it a “A general vectorised if”, which is a good explanation if you’re well-versed in R. But a lot of people aren’t. The best way to describe it is that helps you to avoid nasty nested `if` conditions. You remember when you used Excel2 and you’d want a cell to display something based on an entry in another cell and so you combined `IF()` functions together, well `case_when()` is basically a much nicer and easier to process approach to that. Let’s take our example from the code chunk, we have a set of x and y coordinates and we want to generate the entity given the values of x and y. Let’s also pretend we’re back using Excel (or Sheets, or you spreadsheeting application of choice), so we’ve got a table with our x and y values in columns A and B and we’ve written a formula to provide the entity in column C.\n\n``````| | A | B | C |\n|:--|----:|----:|:--------------|\n| 1 | y | x | entity |\n| 2 | 104 | 36 | retail_recr |\n| 3 | 104 | 210 | grocery_pharm |\n| 4 | 104 | 384 | parks |\n| 5 | 242 | 36 | transit |\n| 6 | 242 | 210 | workplace |\n| 7 | 242 | 284 | residential |\n``````\n\nThe simplest way to write the formula is as follows:\n\n=IF(A2=104,IF(B2=36,“retail_recr”,IF(B2=210,“grocery_pharm”,“parks”)), IF(B2=36,“transit”,IF(B2=210,“workplace”,“residential)))\n\nThere’s a fair potential for error here. Firstly, it only does one test of y, so if the value in A6 was 268 it would still continue to process the second-half of these nested IF statements, the same is also the case for the final set of x values, no matter the value in B4 you’ll get parks if A4 is 104, and in B7 no matter what values you have in either A7 or B7 you’ll get residential.\n\nWe could write this set of IF statements that is more explicit, and will return `FALSE` in any cell where there are incorrect values:\n\n=IF(AND(A2=104,B2=36),“retail_recr”,IF(AND(A2=104,B2=210),“grocery_pharm”, IF(AND(A2=104,B2=384),“parks”,IF(AND(A2=242,B2=36),“transit”, IF(AND(A2=242,B2=210),“workplace”,IF(AND(A2=242,B2=384),“residential”))))))\n\nThis does the job and is error-proof in the sense that it only returns valid results for valid inputs. However, if you’re anything like me then your Excel formulae are lazy and so you’d have written the former. Plus you’re not protected against copy/paste, transposition or overwriting errors.\n\n## case_when() to the rescue\n\nBase R has the `ifelse()` function, and `{dplyr}` has a stricter interpretation `if_else()`, we could nest a set of calls to either of these just as we do in Excel (`ifelse(x==1,\"a\", ifelse(x==2, \"b\", \"c\"))`). But we can avoid messy code by using `dplyr::case_when()`.\n\nIt turns out I use `case_when()` a lot in my code. Some 92 of the 211 `.R` and `.Rmd` files in my live R project folders contain at least one call to `case_when`, a whopping 43.6% of my current scripts! Before today I knew that I used it fairly regularly, but if you’d have asked me to guess I’d probably have said that maybe a quarter or a third of my code used it, not almost half!!\n\nSo what makes `case_when()` so great? Well it’s a very handy function for handling multiple conditions. Let’s take a subset of the code from above, we’ve got a tibble that contains the variables x and y, we’re filtering it to three specific values of y and then creating a new variable called entity using `case_when()`.\n\n``````subnational_datapoints <- subnational_data %>%\nfilter(y == 36 | y == 104 | y == 242) %>%\nmutate(\nentity = case_when(\ny == 36 ~ \"location\",\ny == 104 & x == 36 ~ \"retail_recr\",\ny == 104 & x == 210 ~ \"grocery_pharm\",\ny == 104 & x == 384 ~ \"parks\",\ny == 242 & x == 36 ~ \"transit\",\ny == 242 & x == 210 ~ \"workplace\",\ny == 242 & x == 384 ~ \"residential\")\n``````\n\nWhile you can use `case_when()` working outside of data.frame operations that’s not something I’ve done and I think more conventional nested `base::if()` statements are better, `case_when()` is designed for working with its `{dplyr}` siblings, especially `dplyr::mutate()`. The arguments to `case_when()` are a set of formula pairs, which is where our friend the tilde (`~`) from my colleague’s question comes in. R’s formula notation is something you’ll mainly come across when constructing a linear model, e.g. `y ~ x + z`, but it’s not restricted to its use in model building. A formula is a type of object that R explicitly understands to be a pair of items, a left-hand side (LHS) and a right-hand side (RHS) split by a tilde, which we can generalise this as `LHS ~ RHS`. This two-part nature of a formula pair is what `case_when()` is exploiting.\n\n`case_when()` uses the LHS part as the condition to test for and the RHS as the value to return, it works in sequence through the conditions and when it obtains a `TRUE` result it will return the value associated with that condition. Our first pair, `y == 36 ~ \"location`, tests if y is equal to 36 and if so then it returns “location”. The next pair, `y == 104 & x == 36 ~ \"retail_recr\"`, tests if y is equal to 104 AND x is equal to 36. This code replicates the second Excel formula written above, but I certainly find it easier to read and follow.\n\nThere are a couple of important things to remember when using `case_when()`. Firstly, it’s crucial that the all of the items in your RHS are of the same type, i.e. they’re all character strings, or they’re all numerics, if they’re not the same type then the call will fail. Secondly, if a row in your data doesn’t match any of the conditions then it’ll get an `NA` value, but what if you wanted to have a catch-all value for all other cases, well that’s pretty simple. `case_when()` is testing conditions and looking for a TRUE result, so all you need to do is add a final line of the form `TRUE ~ value`\n\n``````df <- tibble(val = c(1, 2, 1, 3, 4, 5)) %>%\nmutate(new_val = case_when(\nval == 1 ~ \"One\",\nval > 3 ~ \"More than three\",\nTRUE ~ \"Other\"\n))\n``````\n\nIn this code we look at val and if it’s equal to 1 or greater than 3 then we’ll get the response we’ve specified, but in all other cases we’ll get the value “Other”.\n\n``````> df\n\n| val|new_val |\n|---:|:----------------|\n| 1 | One |\n| 2 | Other |\n| 1 | One |\n| 3 | Other |\n| 4 | More than three |\n| 5 | More than three |\n\n``````\n\nSo when’s the case for `case_when()`? Almost half the time if my code is anything to go by.\n\n1. This opens up a segue to the best R prank ever, a tweet that has since been deleted: Today we were learning R Coding and I made a function so that whenever someone typed a ~ it automatically added “Swinton” after it. Three hours it took my colleagues to get it, three hours. - Joe (@raptorbaitjoe) May 29, 2019. ↩︎\n\n2. But not like Kelly Rowland:\n\n↩︎" ]
[ null, "https://lapsedgeographer.london/img/post/2020-04-25-tilda-swinton.gif", null ]
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http://www.expertsmind.com/questions/portion-and-variation-301141728.aspx
[ "## PORTION AND VARIATION, Algebra\n\nAssignment Help:\nSUPPOSE Y IS DIRECTLY PROPORTIONAL TO X AND THAT Y = 35 WHEN X = 5\nFIND THE CONSTANT OF PROPORTIONALITY K\nK=\n\n#### Elimination with 2 problems., {5x-4y=-21} {-2x+4y=18}\n\n{5x-4y=-21} {-2x+4y=18}\n\n(-5,-2) y=5/2x+3\n\n#### Word problems, Five star coffee charges a 40 dollar start up fee and then 8...\n\nFive star coffee charges a 40 dollar start up fee and then 8 dollars per month. Custom coffee company charges a 25 dollar start up fee and 12.50 dollars per month. In which time fr\n\n#### Eponents, How do you write 14x14x14x14x14x14x14= in exponential form?\n\nHow do you write 14x14x14x14x14x14x14= in exponential form?\n\n#### Algebra 2, has a y-intercept of 5 and a slope of 2/3. solve for the standa...\n\nhas a y-intercept of 5 and a slope of 2/3. solve for the standard equation\n\n#### College algebra, how to use the factor theorem\n\nhow to use the factor theorem\n\n#### Areas, does total surface area mean total exposed area\n\ndoes total surface area mean total exposed area\n\n#### Identies, (x+a) (x+b) (x+c) =\n\n(x+a) (x+b) (x+c) =\n\n#### Solving quadratic functions, Sum of two numbers is 10 and their multipicati...\n\nSum of two numbers is 10 and their multipication is 21,find the two numbers. x^2 -10x +21=0\n\n#### Determinant, Power cofactor theorem\n\nPower cofactor theorem", null, "", null, "" ]
[ null, "http://www.expertsmind.com/questions/CaptchaImage.axd", null, "http://www.expertsmind.com/prostyles/images/3.png", null ]
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https://www.bankersadda.com/ibps-rrb-mains-reasoning-quiz-21-2
[ "IBPS RRB PO/Clerk Mains Reasoning Quiz: 21st September 2019\n\nDo you follow a proper plan or strategy for IBPS RRB Mains 2019? Are you aiming IBPS RRB 2019 this time? If yes, then this is the section which can help you to do wonders if practiced well. A good attempt with a mix of accuracy can help you fetch good marks. The reasoning is a game of wits and mind. It is all about logics that a question may have. Speed and accuracy are what that matters the most in this section. The only way to achieve an ambitious goal is by practicing only. So, attempt the quiz of Reasoning ability that inculcates the important questions from the important topics. Do not miss out to practice the Reasoning Ability Quiz that is being provided on Bankersadda.\n\nDirections (1-5): Study the given information and answer the questions:\nWhen a number arrangement machine is given an input line of numbers, it arranges them following a particular rule. The following is an illustration of an input and its rearrangement.", null, "Step IV, is the last step of the above arrangement as the intended arrangement is obtained.\nAs per the rules followed in the given steps find out the appropriate steps for the given input:\n\nInput: 61  22  15  52  71  36  18  11\n\nQ1. What is the sum of the numbers at both the ends in step III of the given arrangement?\n(a) 36\n(b) 33\n(c) 60\n(d) 123\n(e) None of the above\n\nQ2. Which element is 2nd to the right of the one which is 4th from the left end in step I?\n(a) 22\n(b) 44\n(c) 53\n(d) 72\n(e) None of these\n\nQ3. What is the products of the numbers which is 2nd from the right end and 2nd from the left end in final step of the given arrangement?\n(a) 180\n(b) 72\n(c) 13.5\n(d) 180\n(e) None of these\n\nQ4. Which element is 2nd to the left of the one which is 4th from the left end in step II?\n(a) 22\n(b) 44\n(c) 89\n(d) 72\n(e) None of these\n\nQ5. What is the sum of the numbers at both the ends in step IV of the given arrangement?\n(a) 36\n(b) 63\n(c) 60\n(d) 16.5\n(e) None of the above\n\nSolutions (1-5):\nSol. In this input output question numbers are arranged following different rules in each step. Let us understand the logic behind it- In each step the numbers are arranged\nIn step 1: the numbers are multiplied with 3 and then 2 alternatively.\nStep 2: the first two numbers are subtracted, then 2nd and 3rd number are added and alternatively process is repeated till the last number.\nStep 3: Multiplication of the digits of the resultant in the previous step.\nStep4: The numbers in the previous step is divided by 2.", null, "S1. Ans(b)\nS2. Ans(d)\nS3. Ans(e)\nS4. Ans (c)\nS5. Ans (d)\n\nDirections (6-10): In each of the questions given below, a group of digits is given followed by four combinations of letters/ symbols numbered (a), (b), (c) and (d). You have to find out which of the four combinations correctly represents the group of digits based on the letter/ symbol codes and the conditions given below. If none of the four combinations represents the group of digits correctly, give (e) ‘None of these’ as the answer.", null, "Condition for coding the group digits:\n(i) If the first digit is odd and last digit is even, the codes for the first and the last digits are to be interchanged.\n(ii) If the first as well as the last digit is even, both are to be coded by the code for the last digit.\n(iii) If the first as well as the last digit is odd, both are to be coded as ‘X’.\n\nQ6. 285961\n(a) @PD=S\n(b) @∆D=S\n(c) @PV=S\n(d) @PD=SV\n(e) None of these\n\nS6. (a)\nSol. None of the conditions are applied.\n\nQ7. 972486\n(a) =∆@VPS\n(b) S∆@VP=\n(c) SD@VP=\n(d) SA@P=\n(e) None of these\n\nS7. (b)\nSol. Condition (i) Applies.\n\nQ8. 834762\n(a) PMV∆SP\n(b) PMV∆S@\n(c) @MVA∆SP\n(d) @MV∆S@\n(e) None of these\n\nS8. (d)\nSol. Condition (ii) applies\n\nQ9. 785291\n(a) ∆PD@X\n(b) ∆PD@=X\n(c) XPD@=X\n(d) XPD@=\n(e) None of these\n\nS9. (c)\nSol. Condition (iii) applies.\n\nQ10. 748956\n(a) AVP=DS\n(b) SVP=DS\n(c) ∆VP=D∆\n(d) SP=D∆\n(e) None of these\n\nS10. (e)\nSol. Condition (i) applies.\n\nDirections (11-15): Study the following information carefully and answer the given questions.\nA, B, C, D, W, X, Y and Z are eight friends sitting around a square table, two on each side. All of them are facing away from the centre and each is opposite to another. They all live on different floors of a two storey building with topmost floor being numbered as second and the bottommost floor being numbered as first. There are three female members and no two females are seated next to one another. Four of them live on each floor. X sits between D and Z. One of the male members whose immediate neighbours are also males lives on the topmost floor. Y is a female member who sits second to the left of X. Z is not a female member but sits opposite A, who is a female. C sits third to the left of W and is not male member. Only one female lives on the floor below the floor on which D lives.The person who sit opposite to B lives on second floor. A does not lives on topmost floor.C and D are immediate neighbours.\n\nQ11. Who among the following is not a male member?\n(a) W.\n(b) X\n(c) Y\n(d) B\n(e) None of these\n\nQ12. Which of the following statements is true about W and X?\n(a) Both are opposite to each other\n(b) They do not live on the same floor.\n(c) W is female, but X is a male\n(d) Both are females\n(e) None of these\n\nQ13. Which of the following groups includes only females?\n(a) YAW\n(b) ACB\n(c) XYZ\n(d) ACY.\n(e) None of these\n\nQ14. Who among the following is sitting between B and the female whose immediate neighbours live below her, when counted from the right of B?\n(a) A\n(b) C and D\n(c) C\n(d) A and W.\n(e) None of these\n\nQ15. Who amongst the following pair of females lives on the same floor?\n(a) DY\n(b) AY\n(c) AC\n(d) AB\n(e) None of these.\n\nSolutions (11-15):", null, "S11. Ans.(c)\nS12. Ans.(b)\nS13. Ans.(d)\nS14. Ans.(d)\nS15. Ans.(e)\n\nIf you are preparing for Bank exams, then you can also check out a video for Reasoning below:\n\nYou may also like to Read:" ]
[ null, "https://www.bankersadda.com/wp-content/uploads/2019/09/20150723/1.2.png", null, "https://www.bankersadda.com/wp-content/uploads/2019/09/20150725/ans01.png", null, "https://www.bankersadda.com/wp-content/uploads/2019/09/20151214/2221.png", null, "https://www.bankersadda.com/wp-content/uploads/2019/09/20151435/1-to5-300x235.png", null ]
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https://linearalgebras.com/solution-abstract-algebra-exercise-5-1-9.html
[ "If you find any mistakes, please make a comment! Thank you.\n\n## Exhibit symmetric group as a subgroup of a general linear group\n\nSolution to Abstract Algebra by Dummit & Foote 3rd edition Chapter 5.1 Exercise 5.1.9\n\nSolution: We know that if $F$ is a finite field then $\\mathsf{Aut}(F^n) \\cong GL_n(F)$. This isomorphism $\\zeta$ can be defined as follows: given $\\theta \\in \\mathsf{Aut}(F^n)$, $\\zeta(\\theta)$ is the matrix in $GL_n(F)$ whose $i$-th row is precisely $\\theta(e_i)$. (In particular, $\\zeta$ is not canonical since it depends on the choice and order of a basis; here we choose the standard basis.) In Exercises 5.1.7 and 5.1.8, we found an injective group homomorphism $\\psi : S_n \\rightarrow \\mathsf{Aut}(F^n)$. Combining these results we have an injective group homomorphism $\\Psi : S_n \\rightarrow GL_n(F)$, computed as follows: $$\\Psi(\\pi) = [e_{\\pi(1)}\\ \\cdots\\ e_{\\pi(n)}]^T.$$ We can see that each $\\Psi(\\pi)$ is obtained from the identity matrix by permuting the rows, so that each has a single 1 in each row and column and 0 in all other entries. Thus $S_n$ is identified with a subgroup of $GL_n(F)$ consisting of permutation matrices; moreover, by counting we see that all permutation matrices are represented this way.\n\nThus we have proven the result for a finite field $F$; we began with this case because it was shown previously that $\\mathsf{Aut}(F^n) \\cong GL_n(F)$. If this is true for arbitrary fields then the same proof carries over to arbitrary fields $F$; however this will not be proven until later in the text. In the meantime we can convince ourselves that the result holds over arbitrary fields $F$ by noting that, in computing the product of two permutation matrices, we never deal with numbers other than 1 or 0. In particular, the question of whether or not $k = 0$ in $F$ for some integer $k$ never arises, so that all computations hold for arbitrary fields (in which $1 \\neq 0$) and in fact the set of permutation matrices over any $F$ is closed under matrix multiplication." ]
[ null ]
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https://jwcn-eurasipjournals.springeropen.com/articles/10.1186/s13638-021-02017-y
[ "# A four-quadrant mobility model-based routing protocol for post-earthquake emergency communication network\n\n## Abstract\n\nEmergency communication network (ECN) is the important infrastructure to acquire the real-time information after disaster, which is essential for rescue task. However, the existed routing protocols seldom consider the uneven distribution of rescue area so that cannot satisfy the ECN’s requirement and the quality of rescue needs improvement. In this paper, we avoided the traditional linear propulsion from rim to core and proposed a novel four-quadrant mobility model (FQMM), which makes rescuers arrive the most intensity core area first, which can improve the rescue quality. Then, a FQMM-based protocol (FQMMBP) for ECN is designed, which aims to improve the performance of ECN in terms of package delivery rate (PDR) and end-to-end delay. Finally, we set up a virtual earthquake scenario to simulate our proposed protocol in NS2. Results show that the proposed protocol outperforms to the three compared routing protocols, i.e., AODV, DSDV and DSR. The average FDR is improved by 16.31%, and the average reduction of delay is 64.45%, which shows our proposed scheme’s advantages in quality of rescue.\n\n## Introduction\n\n### Background\n\nAfter the devastating earthquake, emergency rescue crews are needed to rush to the scene for rescue as soon as possible. Post-earthquake emergency rescue includes material dispatching, team dispatching, personnel evacuation, post-disaster reconstruction, etc. Due to the different catastrophic degrees caused by the earthquake in different affected areas, the rescue urgency degrees (RUDs) are different for those areas. Seismic intensity, which refers to the catastrophic degree of earthquake impact on the ground or artificial buildings in a certain area , can be used to express the rescue urgencies of disaster areas. The seismic intensity can be obtained in the following ways : In the area where seismic observation equipment is intensively deployed, we can directly obtain the intensity distribution map of the instrument; In the area where seismic observation equipment is sparsely deployed, we can obtain the intensity distribution map of the instrument by grid interpolation; In the area where seismic observation equipment is rarely deployed, we can obtain the intensity distribution map either by historical seismic data statistics or earthquake simulations. Therefore, it is feasible to carry out emergency rescue according to the emergency degree of the disaster area reflected by seismic intensity. When seismic intensity is obtained, the degree of emergency rescue can be determined by the value of seismic intensity, e.g., after an earthquake, the degree of emergency rescue in high intensity areas is higher than that in low intensity areas, while in the region with equal seismic intensity, the degrees of emergency rescue are equal.\n\nAfter the devastating earthquake, the communication infrastructure and power system in the disaster area will suffer from different degrees of damage, resulting in communication system paralysis. Rescue crews have to quickly establish ECN in the disaster area to meet the communication demands both inside and outside the affected area. Unlike traditional self-organized network, the nodes in ECN have the characteristics of energy limitation, high energy consumption in communication and low energy consumption in data calculation. In recent years, the research on ECN routing protocols has been widely concerned by scholars.\n\n### Relate works\n\nIn the existing research on the mobility model of emergency rescue crews, the traditional task assignment model, e.g., dynamic multi-stage allocation model , optimal scheduling strategy and multi-objective optimization model , usually focuses on processing of rescue time. Song Ye et al. established an optimization model for the assignment of earthquake emergency rescue teams, aiming at improving rescue efficiency and maximizing the satisfaction of rescue time. Yuan Jinsha et al. considered demand urgency and satisfaction and constructed a dispatch model of multi-reserve point and multi-disaster point in two-ant group algorithm, which provides a reference to the disaster relief strategy for decision makers. Md. Ibrahim Talukdar et al. investigated shortest path map based (SPMB), random way point (RWP) and random walk (RW) models to measure the efficiency of mobility models. Cao Cejun et al. formulated a bi-level integer programming model (BIPM) to minimize total weighted travel time at the upper level and to maximize total weighted survivors’ perceived satisfaction at the lower level, and then, a case study from Wenchuan earthquake was presented to illustrate the proposed model and solution strategies. Song Yinghua et al. proposed a satisfaction function to measure the psychology of the victims. The emergency material distribution plan optimization model (EMDPOM) was established based on the satisfaction of the victims and the minimum total system cost, combined with the location problem of the emergency distribution center in the two-level distribution network. The model was solved by a real-coded genetic algorithm, and the data in the Wenchuan earthquake were taken as an example to verify the validity and feasibility of the model and algorithm. However, the above researches studied the mobility models from the perspective of rescue capacity and relief materials repertory, without considering seismic intensity in disaster areas. Furthermore, the above researches did not reasonably allocate the communication and rescue nodes based on the characteristics of ECN. Comparisons of traditional mobility models are listed in Table 1.\n\nIn this paper, we first propose a seismic intensity-based four-quadrant rescue mobility model (FQMM), and then, we simulate earthquake scene and propose a FQMM-based protocol for ECN. The main contributions of this paper are as follows.\n\nFirstly, this paper gives the topology of ECN for the scene of post-earthquake emergency rescue. In the proposed post-earthquake ECN, both fixed nodes (e.g., emergency communication vehicles, base stations) and mobile nodes (e.g., Unmanned Aerial Vehicles, portable communication devices) are included. The characteristics of ECN are frequent link breakage, inconsistencies in data rates, incompatibility of resources, the temporary unavailability of needed resources and communications links.\n\nSecondly, the FQMM mobility model for post-earthquake emergency rescue crews is proposed, which based on the topological structure of ECN. With FQMM, first responders not only can carry out grid search in disaster areas to ensure the effectiveness of rescue without missing any victims, but also can carry out radiation diffusion rescue from epicenter to ensure the worst affected areas are on top of the rescue list.\n\nFinally, this paper proposed a FQMM-based protocol (FQMMBP) for ECN. According to the location of the mobile node, FQMMBP predicts hop counts from source to destination, so that it can update the routing table in real time by selecting the optimal next-hop node, thus improving package delivery rate, reducing end-to-end delay.\n\nAlthough the routing protocol proposed in this paper only considers the factor of seismic intensity, the designed mobility model is universal and can be applied to other natural disasters, such as floods, fires, typhoons and so on.\n\n### Problem description\n\nThis paper aims to solve the problems of mobility, communication and rescue task arrangement in the earthquake rescue scenario. The detailed descriptions are as follows:\n\n1. 1.\n\nMobility problem: Based on the common sense of earthquake, the core, rather than the rim of earthquake, has more important emergency to rescue. Therefore, we intend to solve the mobility problem that realizes more rescuers reach the core of earthquake with the minimal traveling time.\n\n2. 2.\n\nRescue task arrangement problem: When the rescuers’ mobility is finished, the paper will solve the problem of rescue task arrangement problem. There is more than one rescuers in each divided zone. And what is the optimal number of arranged rescuers for each zones should be modeled and solved via algorithms.\n\n3. 3.\n\nCommunication problem: Due to the mobility of rescuers, the paper should consider the real-time connectivity communication among the rescue vehicle and rescuers. The proposed routing protocols should guarantee the communication link from rescuers to rescue vehicles during the process of mobility and rescue task, which based on relay connective probability.\n\nOverall, this paper mainly solves the problems defined as above.\n\n## Method\n\n### Scenario and assumption\n\nIn the post-earthquake ECN, the emergency communication vehicle (ECV) and base station (BS) are regarded as fixed nodes, while rescue crews equipped with individual communication equipment are regarded as mobile nodes. Unmanned aerial vehicle (UAV) can be used not only as a fixed node for relaying communication, but also as a mobile node for collecting disaster information. A typical post-earthquake ECN is shown in Fig. 1.\n\nAfter the earthquake, rescuers equipped with individual communication equipment were scattered to various disaster areas for rescue. Each communication node in ECN, such as ECV, UAV and portable individual equipment (PIE), has its own communication coverage. Since ECV cannot reach the center of the disaster area due to traffic control, it starts at the midpoint of length, where rescue crews start and stop rescue missions, as shown in Fig. 1. In the rescue process, ECV can provide real-time communication for rescue crews. Due to its mobility and air superiority, UAV can go deep into the disaster area and play the role of communication relay or disaster information collector. However, due to the limited battery capacity and high costs of UAV, we are unable to deploy a large number of UAVs in the disaster area. Rescuers can enter the disaster area with PIE, however, due to the limited coverage, high cost and large energy consumption, PIE cannot play the role of hot spot or relay for a long time. Therefore, how to cooperate with the communication equipment in ECN to improve the communication efficiency is one of the key issues that need to be paid attention to in post-earthquake emergency rescue. Therefore, heterogeneous communication devices in ECN need to cooperate with each other through appropriate routing protocols to improve communication efficiency.\n\nAfter the earthquake, we use rectangle to represent the earthquake influence range. The length and width of the rectangle are ai and bi (i = 1…n), respectively. Before studying the mobile model, we propose the following assumptions for simulation calculation:\n\n### Assumption 1\n\nIn regions with equal seismic intensity, the values of RUD are equal.\n\n### Assumption 2\n\nEach rescuer has got equal rescue efficiency.\n\n### Assumption 3\n\nRescue tasks in low RUD areas are at high priority. Rescue crews cannot start rescue tasks in higher RUD areas until rescue tasks in low RUD areas are completed.\n\n### Assumption 4\n\nRandom waypoint mobility model for rescuers is used in each disaster area. In order to ensure the fairness and rationality of the rescue sequence, the rescuers are approximately evenly distributed in disaster areas with equal RUD values.\n\n### Assumption 5\n\nLocations and number of rescue tasks are known by rescue crews in advance, since the Global Positioning System (GPS) is used.\n\nBased on seismic intensity and number of rescuers in disaster areas, this paper proposes a four-quadrant mobility model for rescue crews, which improves the rescue efficiency and ensures the safety of the victims' lives and properties. The proposed mobility model mainly includes two aspects: Firstly, the most seriously affected areas are rescued as soon as possible; secondly, the rescue crews are allocated to each appropriate quadrant, which improves the rescue efficiency and shortens the time of rescue.\n\nOur proposed protocols and algorithms are deployed into the rescuers, ECV and UAVs, which use the WSN technologies as platform and infrastructures. Considering the damage of base station and internet on the earthquake scenario, we use ad hoc networks to realize our proposed communication. The ECV is regarded as a tiny server, which is stopping at the rim of the damage area and calculating the RUD and mobility results for each rescuers. And all the rescuers’ mobile information and rescue task status will be send back to the ECV. We assume that every rescuer is implemented an on-broad unit (OBU), which uses WSN technologies. The rescuers can share their basic information, including the real-time position, the task procedure, the remaining energy and the mobile speed, etc. The OBU can realize the rescuers’ communication to other rescuers, ECV and UAVs. As the temporal relay, UAV will improve the communication performance when one of rescuer on the communication link is broken or exhaust energy. UAV will relay the broken communication link and re-arrange the task that exhausted rescuer to other working rescuers.\n\nBased on the overall idea, our proposed algorithms can be implemented by current technologies and can realize effective and real-time communication among rescuers, ECV, and UAVs.\n\n### Disaster area division\n\nDisaster area is divided into n layers according to seismic intensity. The epicenter is located in the first layer (RUD = 1) where length is a1 and width is b1. We define RUD of each region as 1, 2, …, n. Epicenter located at the center of rectangle where RUD = 1. All areas are rectangular rings, except for the one with RUD = 1, which is a rectangle. ECV parks at the midpoint of the length outside the rectangular ring with RUD = n, as shown in Fig. 2.\n\nThe length and width of each disaster area are marked as {RUD|ai, bi}; then, each region can be expressed as follows:\n\n\\begin{aligned} {\\text{Region 1:}} & \\;\\left\\{ {{\\text{RUD}} = 1|a_{1} ,b_{1} } \\right\\}; \\\\ {\\text{Region 2:}} & \\;\\left\\{ {{\\text{RUD}} = 2|a_{2} = a_{1} + \\frac{{a_{1} }}{2} + \\frac{{a_{1} }}{2} = 2a_{1} ,b_{2} = b_{1} + \\frac{{b_{1} }}{2} + \\frac{{b_{1} }}{2} = 2b_{1} } \\right\\}; \\\\ {\\text{Region 3:}} & \\;\\left\\{ {{\\text{RUD}} = 3|a_{3} = a_{2} + \\frac{{a_{2} }}{2} + \\frac{{a_{2} }}{2} = 2a_{2} = 4a_{1} ,b_{3} = b_{2} + \\frac{{b_{2} }}{2} + \\frac{{b_{2} }}{2} = 2b_{2} = 4b_{1} } \\right\\}; \\\\ \\end{aligned}\n\nAccording to mathematical induction, we can get:\n\nRegion i: {RUD = 3|$$a_{i} = 2a_{{i - 1}} = 2^{{i - 1}} a_{1}$$,$${\\text{~}}b_{i} = 2b_{{i - 1}} = 2^{{i - 1}} b_{1}$$}.\n\n### FQMM model\n\nIn this section, a four-quadrant rescuers allocation scheme is proposed. Based on this scheme, a four-quadrant mobility model (FQMM) for rescuers is proposed. The objectives of the allocation scheme are: (1) carrying out grid search in disaster areas to ensure the effectiveness of rescue without missing any victims; (2) carrying out radiation diffusion rescue from epicenter to ensure the worst affected areas are on top of the rescue list. From the disaster area division discussed in Sect. 3.2, it can be seen that when RUD = 1, the region is a rectangle, and when RUD > 1, the regions are rectangular rings. Therefore, the discussion of rescuers allocation scheme includes two aspects: RUD = 1 and RUD > 1. The core ideas of rescuers allocation scheme for regions with RUD > 1 are similar, though the lengths of the regions are different.\n\n#### Mobility model for rescuers in regions with RUD = 1\n\nFive steps (Step 1 to Step 5) are included in the mobility model of rescuers for regions with RUD = 1.\n\nStep 1 As shown in Fig. 3, we establish a coordinate system which takes the epicenter as the origin. X and Y axes are the length and width of the rectangle. Four quadrants, QID = 1, QID = 2, QID = 3 and QID = 4 are counter-clockwise defined according to the definition of coordinate system. Each quadrant is described as follows:\n\n\\begin{aligned} & {\\text{Quadrant 1:}}\\;\\left\\{ {{\\text{RUD}} = 1|{\\text{QID}} = 1} \\right\\};\\;{\\text{ Quadrant 2:}}\\;\\left\\{ {{\\text{RUD}} = 1|{\\text{QID}} = 2} \\right\\}; \\\\ & {\\text{Quadrant 3:}}\\;\\left\\{ {{\\text{RUD}} = 1|{\\text{QID}} = 3} \\right\\};\\;{\\text{Quadrant 4:}}\\;\\left\\{ {{\\text{RUD}} = 1|{\\text{QID}} = 4} \\right\\}. \\\\ \\end{aligned}\n\nSuppose that the total number of rescuers is M, and the total number of rescue tasks in the area with RUD = 1 is Q, which is, $${\\text{RT}}_{{{\\text{QL}} = 0}}^{{{\\text{RUD}} = 1}} = Q$$, where $${\\text{RT}}_{{{\\text{QL}} = 0}}^{{{\\text{RUD}} = 1}}$$ stands for total number of rescue tasks to be completed in disaster areas with RUD = 1 and QL = 0. The number of tasks to be performed in each quadrant is described as: $${\\text{RT}}_{{{\\text{QID}} = k}}^{{{\\text{RUD}} = 1}} = Q/4,\\quad \\left( {k = 1,2,3,4} \\right)$$, where $${\\text{RT}}_{{{\\text{QID}} = k}}^{{{\\text{RUD}} = 1}}$$ stands for the number of tasks to be performed in disaster areas with RUD = 1 and QID = k (k = 1, 2, 3, 4). Rescuers are evenly allocated to four quadrants with QL = 1, and the number of rescuers allocated (NRA) in each quadrant can be discussed from four cases (Case 1 to Case 4) as follows:\n\nCase 1: Mod[M/4] = 0, ere function Mod[x/y] is used to return the remainder of the division of two numbers x, y. In this case, a quarter of rescuers are evenly allocated to four quadrants with QL = 1, which is:\n\n$${\\text{NRA}}_{{{\\text{QID}} = k}}^{{{\\text{RUD}} = 1}} = \\frac{M}{4},\\quad \\left( {k = 1,2,3,4} \\right)$$\n\nwhere $${\\text{NRA}}_{{{\\text{QID}} = k}}^{{{\\text{RUD}} = 1}}$$ stands for the number of rescuers allocated in each quadrant with RUD = 1 and QID = k.\n\nCase 2: Mod[M/4] = 1. In this case, one more rescuer is allocated to Quadrant 1, which are:\n\n$${\\text{~NRA}}_{{{\\text{QID}} = 1}}^{{{\\text{RUD}} = 1}} = \\frac{{M - 1}}{4} + 1$$\n\nwhere $${\\text{NRA}}_{{{\\text{QID}} = 1}}^{{{\\text{RUD}} = 1}}$$ stands for the number of rescuers allocated in Quadrant 1 with RUD = 1, and $${\\text{NRA}}_{{{\\text{QID}} \\ne 1}}^{{{\\text{RUD}} = 1}}$$ stands for the number of rescuers allocated in Quadrant 2, Quadrant 3 and Quadrant 4 with RUD = 1.\n\nCase 3: Mod[M/4] = 2. In this case, one more rescuer is allocated to Quadrant 1 and Quadrant 2, respectively, which are:\n\n$${\\text{NRA}}_{{{\\text{QID}} = 1\\;{\\text{or}}\\;2}}^{{{\\text{RUD}} = 1}} = \\frac{{M - 2}}{4} + 1$$\n\nwhere $${\\text{NRA}}_{{{\\text{QID}} = 1\\;{\\text{or}}\\;2}}^{{{\\text{RUD}} = 1}}$$ stands for the number of rescuers allocated in Quadrant 1 and Quadrant 2 with RUD = 1, and $${\\text{NRA}}_{{{\\text{QID}} = 3\\;{\\text{or}}\\;4}}^{{{\\text{RUD}} = 1}}$$ stands for the number of rescuers allocated in Quadrant and Quadrant 4 with RUD = 1.\n\nCase 4: Mod[M/4] = 3. In this case, one more rescuer is allocated to Quadrant 1, Quadrant 2, Quadrant 4 which are:\n\n$${\\text{NRA}}_{{{\\text{QID}} \\ne 3}}^{{{\\text{RUD}} = 1}} = \\frac{{M - 3}}{4} + 1$$\n\nwhere $${\\text{NDN}}_{{{\\text{QID}} \\ne 3}}^{{{\\text{RUD}} = 1}}$$ stands for the number of rescuers allocated in Quadrant 1, Quadrant 2 and Quadrant 4 with RUD = 1, and $${\\text{NRA}}_{{{\\text{QID}} = 3}}^{{{\\text{RUD}} = 1}}$$ stands for the number of rescuers allocated in Quadran3 with RUD = 1.\n\nStep 2 According to the procedure of step 1, each quadrant with RUD = 1 is further divided into four secondary quadrants as shown in Fig. 4. Two elements are included in QID set, which record the quadrant ID from the first level to the last level. For example, {RUD = 1|QID = {1, 3}} in Fig. 4 means region with QL = 1 is located in Quadrant 1 and region with QL = 2 is located in Quadrant 3.\n\nStep 3 For QL = x, replace the value of M with $${\\text{RT}}_{{{\\text{QL}} = x}}^{{{\\text{RUD}} = 1}}$$, repeat Step 1 and calculate the number of rescuers allocated for each quadrant with QL = x.\n\nStep 4 Repeat Step 2 and Step 3, until any of the following Abort conditions is satisfied.\n\nAbort condition 1 There is only one rescuer in each area of quadrant with QL = x, which is: $${\\text{NRA}}_{{{\\text{QID}} = k}}^{{{\\text{RUD}} = 1}} = 1,\\quad \\left( {k = 1,2,3,4} \\right)$$, where $${\\text{NRA}}_{{{\\text{QID}} = k}}^{{{\\text{RUD}} = 1}}$$ is the number of rescuers allocated in Quadrant k (k = 1, 2, 3, 4) with RUD = 1 and QL = x. Under this condition, further quadrant division is no longer meaningful, so the single rescuer will complete all rescue tasks in the region of quadrant with QL = x.\n\nAbort condition 2 The number of rescue tasks is less than that of rescuers assigned in the quadrant with QL = x, which is:$$~{\\text{NRA}}_{{{\\text{QID}} = k}}^{{{\\text{RUD}} = 1}} \\ge {\\text{RT}}_{{{\\text{QID}} = k}}^{{{\\text{RUD}} = 1}} ,\\quad \\left( {k = 1,2,3,4} \\right)$$, where $${\\text{NRA}}_{{{\\text{QID}} = k}}^{{{\\text{RUD}} = 1}}$$ is the number of rescuers allocated in Quadrant k (k = 1, 2, 3, 4) with RUD = 1 and QL = x, and $${\\text{RT}}_{{{\\text{QID}} = k}}^{{{\\text{RUD}} = 1}}$$ is number of rescue tasks in Quadrant k (k = 1, 2, 3, 4) with RUD = 1 and QL = x. Under this condition, each rescuer is assigned to at most one rescue task, and there are unoccupied rescuers in the area.\n\nStep 5 Quadrant division is suspended and rescuers begin to perform their own rescue tasks.\n\nFigure 5 shows the algorithm flow of rescue node allocation in rectangular disaster area when RUD = 1.\n\n#### Mobility model for rescuers in regions with RUD > 1\n\nAccording to the mobility model, rescue crews cannot start rescue tasks in higher RUD areas until rescue tasks in low RUD areas are completed. As shown in Fig. 6, in the gray area with RUD = 1, rescuers are carrying out rescue tasks. When the task is over, rescuers will be assigned to the white rectangular ring area with RUD = 2.\n\nFive steps (Step 6 to Step 10) are included in the mobility model of rescuers for regions with RUD > 1.\n\nStep 6 As shown in Fig. 7, we establish four coordinate systems which takes O1, O2, O3 and O4 as the origins, respectively. Then, the rectangular ring is equally divided into twelve quadrants (the white parts in Fig. 7) by four coordinate systems. In each coordinate system, four quadrants are equal with each other in area, as shown in Fig. 7 where four triangle marks.\n\nStep 7 It can be seen from Step 6 that in each coordinate system, there are three quadrants waiting for allocation of rescuers, and the left quadrant is allocated with rescue crews. As shown in Fig. 7, in the coordinate system with O3 as origin, there are four quadrants where Triangle (1), Triangle (2), Triangle (3) and Triangle (4) marks. Area where Triangle (1) marks is allocated with rescue crews, while areas where Triangle (2), Triangle (3) and Triangle (4) marks are waiting for rescuers allocations. We take coordinate system O3 as an example. Suppose there are M rescuers in Quadrant 1 with RUD = i − 1, which marked with Triangle (1). There are Q rescue tasks in areas with RUD = i, which is: $${\\text{RT}}_{{{\\text{QL}} = 0}}^{{{\\text{RUD}} = i}} = Q$$. Rescuers are evenly allocated to three quadrants [Triangle (2), Triangle (3) and Triangle (4)] with QL = 1, and the number of rescuers allocated (NRA) in each quadrant can be discussed from three cases (Case 5 to Case 7) as follows:\n\nCase 5 Mod[M/3] = 0, where function Mod[x/y] is used to return the remainder of the division of two numbers x, y. In this case, one-third rescuers are evenly allocated to three quadrants, which is:\n\n$${\\text{NRA}}_{{{\\text{QID}} = k}}^{{{\\text{RUD}} = i}} = \\frac{M}{3},\\quad \\left( {k = 2,3,4} \\right)$$\n\nwhere $${\\text{NRA}}_{{{\\text{QID}} = k}}^{{{\\text{RUD}} = i}}$$ stands for the number of rescuers allocated in each quadrant with RUD = i and QID = k.\n\nCase 6 Mod[M/3] = 1. In this case, one more rescuer is allocated to Quadrant 2, which are:\n\n$${\\text{NRA}}_{{{\\text{QID}} = 2}}^{{{\\text{RUD}} = i}} = \\frac{{M - 1}}{3} + 1$$\n\nwhere $${\\text{NRA}}_{{{\\text{QID}} = 2}}^{{{\\text{RUD}} = i}}$$ stands for the number of rescuers allocated in Quadrant 2 with RUD = i, and $${\\text{NRA}}_{{{\\text{QID}} = 3\\;{\\text{or}}\\;4}}^{{{\\text{RUD}} = i}}$$ stands for the number of rescuers allocated in Quadrant 3 and Quadrant 4 with R = i.\n\nCase 7 Mod[M/3] = 2. In this case, one more rescuer is allocated to Quadrant 2 and Quadrant 4, respectively, which are:\n\n$${\\text{NRA}}_{{{\\text{QID}} = 2\\;~{\\text{or}}\\;4}}^{{{\\text{RUD}} = i}} = \\frac{{M - 2}}{3} + 1$$\n\nwhere $${\\text{NRA}}_{{{\\text{QID}} = 2~\\;{\\text{or}}\\;4}}^{{{\\text{RUD}} = i}}$$ stands for number of rescuers allocated in Quadrant 2 and Quadrant 4 with RUD = i, and $${\\text{NRA}}_{{{\\text{QID}} = 3}}^{{{\\text{RUD}} = i}}$$ stands for the number of rescuers allocated in Quadrant 3 with RUD = i.\n\nThe number of rescue tasks in each quadrant with RUD = i is:$${\\text{RT}}_{{{\\text{QID}} = k}}^{{{\\text{RUD}} = i}} = Q/12,\\quad \\left( {k = 2,3,4} \\right)$$, where $${\\text{RT}}_{{{\\text{QID}} = k}}^{{{\\text{RUD}} = i}}$$ stands for the number of rescue tasks in Quadrant k (k = 2, 3, 4) wh RUD = i.\n\nStep 8 Quadrants (QL ≥ 2) division methods are the same as illustrated in Step 6, which take the middle points of quadrants with QL = 1 as the origins and establish the coordinate systems. Repeat Step 7, until any of the following Abort conditions is satisfied.\n\nAbort condition 3 There is only one rescuer in each area of quadrant with QL = x, which is:$${\\text{NRA}}_{{{\\text{QID}} = k}}^{{{\\text{RUD}} = i}} = 1,\\quad \\left( {k = 1,2,3,4} \\right)$$, where $${\\text{NRA}}_{{{\\text{QID}} = k}}^{{{\\text{RUD}} = i}}$$ is the number of rescuers allocated in Quadrant k (k = 1, 2, 3, 4) with RUD = i and QL = x. Under this condition, further quadrant division is no longer meaningful, so the single rescuer will complete all rescue tasks in the region of quadrant with QL = x.\n\nAbort condition 4 The number of rescue tasks is less than that of rescuers assigned in the quadrant with QL = x, which is:$${\\text{NRA}}_{{{\\text{QID}} = k}}^{{{\\text{RUD}} = i}} \\ge {\\text{RT}}_{{{\\text{QID}} = k}}^{{{\\text{RUD}} = i}} ,\\quad \\left( {k = 1,2,3,4} \\right)$$, where $${\\text{NRA}}_{{{\\text{QID}} = k}}^{{{\\text{RUD}} = i}}$$ is the number of rescuers allocated in Quadrant k (k = 1, 2, 3, 4) with RUD = i and QL = x, and $${\\text{RT}}_{{{\\text{QID}} = k}}^{{{\\text{RUD}} = i}}$$ is number of rescue tasks in Quadrant k (k = 1, 2, 3, 4) with RUD = i and QL = x. Under this condition, each rescuer is assigned to at most one rescue task, and there are unoccupied rescuers in the area.\n\nStep 9 Quadrant division is suspended and rescuers begin to perform their own rescue tasks.\n\nFigure 8 shows the algorithm flow of rescue node allocation in rectangular disaster area when RUD > 1.\n\n### FDMM-based routing protocol\n\nIn this section, a proposal of FQMM-based protocol (FQMMBP) for ECN is presented. First, we discuss route discovery and maintenance of FQMMBP, and then we propose FQMMBP protocol by analyzing the communication mechanism of nodes in ECN.\n\nPost-earthquake ECN, which is a hybrid self-organized network, mainly includes ECV, UAV, portable devices carried by rescue crews, base station, satellite, etc. Those communication devices can be classified as fixed communication nodes (such as ECV, base station, UAV used for relay.) and mobile communication nodes (such as portable devices carried by rescue crews, satellite.). As shown in Fig. 1, each portable device is equipped with a wireless communication module. Rescuers can use portable devices to build self-organized networks, through which they can communicate with each other and complete the transmission process of data packets, such as receiving, transmitting and relaying. Rescuers can also communicate with ECV, UAV and other fixed nodes within the coverages of portable devices. As a fixed communication node, ECV can communicate with rescuers and UAVs in the disaster area through downlink, or with the base station, satellite and other communication facilities through uplink. Through the backbone network, commanders can get the information about rescue crews, rescue process and suffering condition in time, so as to make scientific rescue decision and emergency response quickly.\n\nWe define the communication radius of mobile node as R. The mobile node sends data to the emergency communication vehicle in the way of multi-hop. In such communication mode, the source is the mobile node in the disaster area, the destination is ECV at the edge of the disaster area, and the relay nodes are mobile nodes or fixed nodes (UAVs, etc.) in the disaster area. On the premise of ensuring the stability of communication links, we choose relay nodes or links that can minimize the delay for data transmission. Based on the mobility model we discussed in Sect. 3.3, it can be seen that rescue crews are nearly evenly distributed in disaster areas with the same RUD. In order to ensure the stability of the communication link in the rescue process, we place fixed UAVs at the edges of disaster areas with different RUD values; however, they do not undertake the rescue tasks. When all rescue tasks in the disaster area with a specific RUD value are completed, the relay UAV in the area returns to the location of ECV.\n\n#### Routing setup\n\nIn order to reduce the data transmission delay, we select the next hop relay node according to RUD value, QL and QID of each region where mobile nodes are located in. RUD and QL are used to define the level of the region where mobile nodes are located in. QID is a set which includes 1, 2, 3 and 4, and there are QL elements in the set. According to the set, the mobile nodes learn the division of the regions and the quadrants position, so as to select the appropriate next hop mobile node. Therefore, in the process of routing discovery, it is necessary to extend the routing request (RREQ) message, so as to meet the requirements of our proposed communication mechanism of FQMMBP. We improved the traditional RREQ message by adding four new fields: RUD, QL, QID and QW. RUD describes the area where nodes are located and judge whether the area is a rectangle or a rectangular ring as shown in Fig. 2. QL describes the level of the quadrant where nodes are located, and it determines the relative position of the mobile rescue node. QID describes the division of the quadrant where nodes are located and judge its relative position to ECV. When the QL of quadrants are the same, QID helps pick up the mobile rescue node which is closer to ECV. QW calculates weights of nodes according to QL and QID, so as to select the most appropriate next-hop relay node.\n\nIn the process of route discovery, nodes broadcast RREQ messages to send request information to their neighbors. The improved RREQ data format is shown in Table 2.\n\nEach field of RREQ data format is described as follows:\n\n• Type: The length of this field is 8bit, and the type value of RREQ is 1.\n\n• Type: The length of this field is 8bit, and the type value of RREQ is 1.\n\n• Flags: The length of this field is 5 bits, and five identities (J, R, G, D, and U) are included in this field. \"J\" is a joint flag and is generally used for multicast. “R” is for route repairing and is used for multicast transmission. “G” represents the list of nodes around ECV which can be communicated. Flag “G” determines whether the RREQ message can be directly sent to the destination or not. “D” is the reply flag of the destination node which determines whether the destination node is allowed to reply to the received RREQ message or not. “U” is the flag of unknown serial number. U = 1 means that the serial number of the node is unknown.\n\n• Reserved: Reserved field for further improvement of RREQ message.\n\n• Hop Count: This field registers the hop counts that RREQ passes from the source node to the current node.\n\n• RREQ ID: This field is the unique identity of RREQ message.\n\n• RUD: Rescue Urgency Degree, which is inversely proportional to the seismic intensity value in this area.\n\n• QL: Quadrant Level, which is the number of times the quadrant divided.\n\n• QW: Weights of nodes in quadrant.\n\n#### Communication mechanism\n\n$$D_{{{\\text{QID}}}} = \\left\\{ {\\begin{array}{*{20}l} c \\hfill & {{\\text{QID}} = 1\\;{\\text{or}}\\;2} \\hfill \\\\ d \\hfill & {{\\text{QID}} = 3\\;{\\text{or}}\\;4} \\hfill \\\\ \\end{array} } \\right.$$\n(1)\n$$c = \\sqrt {\\left( {\\frac{a}{4}} \\right)^{2} + \\left( {\\frac{b}{4}} \\right)^{2} }$$\n(2)\n$$d = \\sqrt {\\left( {\\frac{a}{4}} \\right)^{2} + \\left( {\\frac{{3b}}{4}} \\right)^{2} }$$\n(3)\n\nwhere a, b are length and width of disaster area, respectively.$$~D_{{PH}}$$ stands for the distance between the center of each quadrant and ECV. c is the distance between the center of Quadrant 1 and ECV, which equals with the distance between the center of Quadrant 2 and ECV. d is the distance between the center of Quadrant 3 and ECV, which equals with the distance between the center of Quadrant 4 and ECV. Take the weight of Quadrant 1 (or Quadrant 2) as the benchmark, which is: $$W_{{{\\text{QID}} = 1\\;{\\text{or}}\\;2}} = 1$$, where $$W_{{{\\text{QID}} = 1\\;{\\text{or}}\\;2}}$$ is the weight of Quadrant 1 or Quadrant 2. The weight of Quadrant 3 (or Quadrant 4) can be expressed as:\n\n$$W_{{{\\text{QID}} = 3\\;{\\text{or}}\\;4}} = \\frac{d}{c} = \\frac{{\\sqrt {a^{2} + 9b^{2} } }}{{\\sqrt {a^{2} + b^{2} } }}$$\n(4)\n\nwhere $$W_{{{\\text{QID}} = 3\\;{\\text{or}}\\;4}}$$ is the weight of Quadrant 3 or Quadrant 4.\n\nDefinition $${\\text{QID}}_{i}$$ is the Quadrant ID with QL = i. For example, $${\\text{QID}}_{3} = 4$$ means Quadrant 4 is with QL = 3. From Eqs. (1) and (4), we can get:\n\n$$D_{{{\\text{QID}}_{{\\text{i}}} }} = \\left\\{ {\\begin{array}{*{20}l} {c_{i} } \\hfill & {{\\text{QID}}_{i} = 1\\;{\\text{or}}\\;2} \\hfill \\\\ {d_{i} } \\hfill & {{\\text{QID}}_{i} = 3\\;{\\text{or}}\\;4} \\hfill \\\\ \\end{array} } \\right.$$\n(5)\n\nwhere $$D_{{{\\text{QID}}_{i} }}$$ stands for the distance between the center of quadrant with QIDi and ECV. ci is distance between the center of Quadrant 1 with QL = i and ECV, which equals with the distance between the center of Quadrant 2 with QL = i and ECV. di is distance between the center of Quadrant 3 with QL = i and ECV, which equals with the distance between the center of Quadrant 4 with QL = i and ECV. From Eq. (4), we can get:\n\n$$W_{{{\\text{QID}}_{i} }} = \\frac{{\\sqrt {a_{i} ^{2} + 9b_{i} ^{2} } }}{{\\sqrt {a_{i} ^{2} + b_{i} ^{2} } }}$$\n(6)\n\nwhere $$W_{{{\\text{QID}}_{i} }}$$ is the weight of node located in Quadrant $${\\text{QID}}_{i}$$. ai, bi are length and width of disaster area with QL = i, respectively. According to the QID and QL, we can get the total weight of each node in quadrants with different levels by:\n\n$${\\text{QW}} = \\mathop \\sum \\limits_{{i = 1}}^{{{\\text{QL}}}} W_{{{\\text{QID}}_{i} }}$$\n(7)\n\nwhere QW is the total weight of one node in quadrants with different levels.$${\\text{~}}W_{{{\\text{QID}}_{i} }}$$ is the weight of node located in Quadrant $${\\text{QID}}_{i}.$$\n\n#### Routing discover and maintain\n\nIn this section, we propose FQMMBP protocol based on FQMM. Routing discovery and routing maintenance process are illustrated in Figs. 10 and 11, respectively.\n\nAs shown in Fig. 10, source node (denoted as NS) checks whether there is a route to ECV. If there is no route to ECV, the node broadcasts RREQ message to its neighbor node (denoted as NN). After receiving the message, NN first checks whether NS is a valid node. If NS is invalid, NN discards the received RREQ message; otherwise, NN updates the routing table and finds out whether there is routing information to ECV. If there is no route to ECV, NN forwards the received RREQ message. Repeat the broadcasting and checking process until one or more valid routes to the ECV are found.\n\nAs shown in Fig. 11, source node (denoted as NS) sends hello packets to the candidate nodes within its coverage by broadcasting and receives RREQ replies from them. NS determines whether there is a RREQ reply sent by ECV by checking whether there is a candidate node with RUD = n and both QL and QID are empty in RREQ reply. If there is an ECV within the coverage of NS, data packets are directly transmitted to ECV which is the destination of the communication, and communication ends. If there is no ECV within the coverage of NS, then NS determines whether there is a RREQ reply sent by fixed UAV by checking whether there is a candidate node with RUD ≠ n and both QL and QID are empty in RREQ reply. If there is an UAV within the coverage of NS, data packets are forwarded to UAV, and UAV forwarded data packets to its neighbor. Repeat this process until ECV is founded, finally communication ends. If there are neither ECVs nor UAVs in the candidate nodes, we compare QW values. If there is only one candidate node corresponding to the maximum QW value, then this node is the optimal next hop node (denoted as NOPT). NOPT becomes a new source node NS, and then it sends the hello packets to repeat the process. If the maximum QW value corresponds to more than one candidate nodes, we should update candidate nodes set by removing nodes whose QW is not the maximum. Moreover, the protocols flowchart adds a judgment module to take the energy consumption into consideration. The model will test the candidate relays’ remaining energy and selects the suitable candidate as best relay to achieve the communication so that the ECNs which remain the fewer energy will be weeded out from the candidate set to guarantee the quality of rescue. After that, we select NOPT from the new set according to the following Cases:\n\nCase 1: There is only one candidate node with the maximum number of Quadrant 1. This node is selected as NOPT.\n\nCase 2: The numbers of Quadrant 1 are equal, while there is only one candidate node with the maximum number of Quadrant 2. This node is selected as NOPT.\n\nCase 3: The numbers of Quadrant 2 are equal, while there is only one candidate node with the maximum number of Quadrant 4. This node is selected as NOPT.\n\nCase 4: The numbers of Quadrant 4 are equal, while there is only one candidate node with the maximum number of Quadrant 3. This node is selected as NOPT.\n\nCase 5: Recalculate the weights of each candidate node according to MCDM-ECP algorithm and then select NOPT.\n\nThe description of FQMMBP routing maintenance process is completed.\n\n### An example\n\nIn order to understand the entire processes of our propose FQMMBP protocol clearly, we give an example to illustrate the main processes including the disaster area division, the main mobility in different RUD, and the relay selection in communication link based on FQMM.\n\nAs shown as Fig. 12a, we give an example of a disaster area of 6 km * 4 km. According to the degree of emergency importance, the 1 km * 2 km core of earthquake will be divided as RUD = 1, which means the mobile rescuers will do the task in this area firstly. Based on the equation forehead Sect. 2.2, we can determine and calculate the length and width of each RUD. Besides RUD = 1 area is rectangular, all the areas of RUD > 1 are rectangular rings. Based on the results of disaster area division, the smaller RUD id will be rescued earlier. For this instance, we divide 3 disaster areas, RUD = 1 will be rescued first. RUD = 2 area is rescued once the tasks in RUD = 1 area are finished. So are the tasks in area of RUD = 3.\n\nAs shown as Fig. 12b, we suppose that 50 rescuers are set into the disaster area. When RUD = 1, the shape of area is rectangular, we first divide the area into four quadrants. According to the formulae in Sect. 2.3.1 50 mod 4 = 2. Therefore, the remaining 2 rescuers is arranged into first and second quadrants. As a result, 50 rescuers are arranged into four quadrants as 13, 13, 12 and 12 rescuers.\n\nAs shown as Fig. 12c, RUD = 2 and RUD = 3 are the same rectangular rings, which the rescuers are focused in one quadrant of the whole sub-quadrants. For the upper 2 sub-quadrants, 13 mod 3 = 1. Therefore, 13 rescuers are arranged to the other 3 quadrants as 5, 4 and 4 rescuers. For the lower 2 sub-quadrants, 12 mod 3 = 0. Therefore, the 12 rescuers mobile into the other 3 quadrants evenly. Attentionally, the rescuers must mobile and begin the next RUD’s tasks after the last RUD’s tasks are finished. And the additional one rescuer is often arranged to the quadrants that are nearly to the central line.\n\nAs shown as Fig. 12d, the communication link is setup from the communication vehicle that located at the rim of disaster area to the rescuers in four quadrants. The core communication link goes throughout the central line of disaster area so that the conditions of all quadrants can be acquired. Based on the formulae in Sect. 2.4, every rescuers as a sequence to identify its own position in the quadrants. The routing protocol will select the suitable relay nodes, as shown as red triangular in the figure to realize the communication to rescue vehicle. As cluster head, these relay nodes will also collect the real-time traffic and rescue information in their own sub-quadrants.\n\n## Results and discussion\n\nIn this section, we present the network level performance evaluation and simulations results of the proposed routing protocol based on NS-2 platform. For comparison, with the respect to each relevant class of routing protocols, we selected the pertinent protocols which are the most widely used (e.g., AODV, DSR, DSDV).\n\n### Experiment setup\n\nThe simulation environment is a 10 km × 10 km earthquake disaster area, in which the epicenter of the most severely affected area is 6 km × 4 km. There are 48 mobile rescue nodes, two fixed UAVs and one ECV. The simulation setup and respective parameters are detailed in Table 3.\n\n### Experiment results\n\nIn this section, four routing protocols, e.g., FQMMBP, AODV, DSDV and DSR, are analyzed from the perspective of package delivery rate (PDR), end-to-end delay (delay) and overhead.\n\n#### Package delivery rate\n\nFigure 13 shows the performance of Package Delivery Rate (PDR) for the different studied protocols. Our proposed FQMMBP achieves the best performance with nearly 91.24% of average PDR while the other three protocols with about 80% of average PDR. This proves that FQMMBP improves communication efficiency and stabilizes communication link status in a long period of time. DSDV achieves the worst performance with 79.23% of average PDR mainly due to the mobility of communication nodes. DSDV is a proactive protocol, and each node maintains a route table. In the emergency rescue scenario, routing tables are dynamic due to the mobility of rescuers, thus causing the cost of maintaining routing tables. The superiority of FQMMBP over other three routing protocols is mainly attributed to FQMM mobility model. In each disaster area under this model, the distribution of rescue mobile nodes tends to be uniform, which increases the probability that each node can find the appropriate next-hop node to forward data and ensures the integrity and reliability of communication link. Packets are likely to be delivered successfully in such distributions of nodes.\n\nCompared to other protocols, our proposed FQMMBP has the best performance in FDR, which is increased to AODV, DSDV and DSR as 15.81%, 17.59% and 15.54%, respectively. Therefore, our proposed routing protocol can improve average FDR as 16.31%. And when the simulation time denotes 450 s, our FQMMBP has the best increments in FDR compared to DSDV, which rises 32.87%. However, the FDR performance of FQMMBP at the beginning of simulation has the fewer advantage, because the rescuers should move into the core area of earthquake rather than the rim of scenario, which the FDR will be weak.\n\nPDR performances of the four compared protocols are almost the same before 150 s, that is because all nodes are concentrated in the central area with RUD = 1 at the beginning of simulation. At this point, the distributions of nodes in the four protocols are similar, and each node has not started to disperse or move to the next area according to the mobility model. As time goes on, the PDR of each protocol decreases gradually. This is because when the nodes complete the rescue tasks in the central area, they will spread to the next peripheral rectangular ring, resulting in the scattered distribution of nodes and the increase in distance between nodes, so the stability of communication link becomes poor. On the contrary, mobile nodes are evenly distributed and close to ECV under our proposed FQMMBP, which improves PDR.\n\n#### End-to-end delay\n\nFigure 14 shows the performance of End-to-end delay (delay) for the different studied protocols. In terms of delay, FQMMBP is significantly lower than other classical routing protocols. During the simulation, the average delay decreased from 1.44 s of DSR to 0.55 s of FQMMBP, which proves the significant improvement of communication performance. This is because in each rescue area, the distribution of mobile nodes tends to be uniform under FQMMBP, which increases the probability that each node can find the next hop node to communicate with. Within the coverage of each relay node, the existence probability of candidate next hop node increases, which shortens the time to find the candidate next hop node as well as the time to transmit data to ECV.\n\nCompared to other protocols, our proposed FQMMBP has the best performance in delay, which is reduced to AODV, DSDV and DSR as 60.20%, 66.56% and 66.58%, respectively. Therefore, our proposed routing protocol can decrease average delay as 64.45%. And when the simulation time denotes 400 s, our FQMMBP has the most decrement in delay compared to DSR, which is 1.459393 s. The reason of the trend of FQMMBP at the beginning is also above the rescuers’ mobility.\n\nPerformances of four compared protocols are almost the same before 150 s. This is because at the beginning of the simulation, all nodes are concentrated in the central rectangular area with the most serious seismic intensity. At this time, the distributions of nodes are relatively centralized and similar, and the nodes have not started to disperse and move according to their respective mobility models, so the performance differences of the four protocols are not significant. As time goes on, the delay of each protocol increases gradually, because all nodes will spread to the next peripheral rectangular ring when the rescue tasks in the central area are completed. The dispersion of nodes leads to the increase in distance between nodes; therefore, it takes longer for relay nodes to find the appropriate next hop nodes as the decreasing numbers of candidate next hop nodes. The distribution of mobile nodes is more uniform under the proposed FQMMBP, which shortens the delay.\n\nIn the period of 200–300 s, the delay increases obviously as the nodes spread to the next area around when they finish the rescue tasks in the central area. The node distribution changes from the original rectangular centralized to the rectangular ring decentralized, so the node distribution density becomes low, which increases the delay. After 300 s, the delay tends to be stable. This is because there is no significant change in the distributions of nodes when they move from the area with RUD = 2 to the areas with RUD > 2. In that case, the delay tends to be stable because of the little impact on FQMMBP.\n\nFigure 15 shows the performance of overhead for the different studied protocols.\n\nIn terms of overhead, FQMMBP is significantly larger than other classical routing protocols. That is because four new fields (RUD, QL, QID and QW) are added to RREQ message. Furthermore, the dispersion of nodes increases the number of forwarding packets, which also contributes to the increase in overhead. The data transmission overhead of AODV, DSDV and DSR protocols increases with time, while FQMMBP meets the same rule from 100 to 300 s. However, the overhead of FQMMBP is decreasing between 300 and 450 s, that is because before 300 s, the rescue nodes are distributed centrally, and they are far away from the destination nodes. Most of the data transmissions rely on fixed relay nodes for forwarding, so the overhead is increasing. In the period of 300–450 s, the distribution of rescue nodes is scattered, and the average distance between rescue nodes and destination nodes is close. In that case, the number of data forwarding is reduced, so the overhead is reduced. After 450 s, the distribution of rescue nodes is more scattered than before, and the data communication between nodes relies on multi-hop forwarding, which makes the data transmission overhead increasing. After 600 s, the overhead of FQMMBP can be controlled within 2.5%.\n\nAs a negative impact of lower delay and higher FDR, our FQMMBP lost some extents of overhead, which is a metric relating to the energy and rescuer’s cost. Compared to other protocols, our proposed FQMMBP has increased overhead to AODV, DSDV and DSR as 49.89%, 46.83% and 43.92%. Our proposed routing protocol has average 46.88% increased overhead. We trade 46.88% overhead loss to gain the benefits of 64.45% delay and 16.31% FDR. It is a remarkable and acceptable results for the specific earthquake scenario, which the victims’ life is much more important than rescuers’ battery. Although the performance of FQMMBP in overhead is not as good as the other three protocols, it is worthwhile for the emergency rescue.\n\n### Discussions\n\n#### Energy consumption\n\nThere is no doubt that energy consumption is an important metric in mobile rescue issue. In our previous work, we have analyzed the energy consumption during the communication and mobility .\n\nOverall in two sides, the energy consumption is taken into consideration in our previous work, and the energy exhaustion problem can be solved in the realistic emergency rescue scenario.\n\n#### Performance metrics selection\n\nWe have two anticipations for our proposed protocols: one is that the ECNs can be set up as soon as possible, which is related to the metric of delay. This metric (delay) can guarantee the rescuers acquire and collect the real-time information about the disaster and victim with short time. The other is the communication link can maintain as well as possible, which means the ECNs can acquire the sufficient information and related to the metric of FDR. For this work, all selected performance metrics (delay and FDR) are useful for the victims in the earthquake. For general scenario, energy consumption is a vital metric for communication modules. The lower energy consumption, the longer modules can remain. However, in this specific earthquake scenario, the lives of victims (human) are much more important than the modules’ (electric rescuers). We should guarantee the human being’s life with the fast (lower delay) and successful (high FDR) information, not one communication module’s duration in the communication link. The life of human is more vital and indispensable than the life of communication module. Actually, in order to save victims lives, we scarify ECN’s energy as the addition cost. This is our primary reason to select the delay and FDR as metrics, rather than energy consumption. Moreover, based on our revision, which took the energy consumption into consideration, the protocol’s selection results are not changed.\n\n#### Routing protocols selection\n\nIn our previous work’s simulation results, the selection policy has the lowest energy consumption. Therefore, we improve the selection policy by considering the rescuers’ mobility model and update it as FQMMBP. The idea of disaster area’s division, and the four-quadrant mobility model are our main contribution about the proposed routing protocol. This protocol has two advantages: one is that the protocol inherits the previous lower energy consumption advantage. The other is that the protocol can guarantee the quality of rescue in earthquake because of the disaster area division and rescuers’ mobility model.\n\nThe other routing protocols cannot work well in the earthquake emergency scenario. DSDV is adapted from the conventional Routing Information Protocol (RIP) to ad hoc networks routing. It adds a new attribute, sequence number, to each route table entry of the conventional RIP. Using the newly added sequence number, the mobile nodes can distinguish stale route information from the new and thus prevent the formation of routing loops. DSR is set up with the route discovery and route maintenance mechanism to make sure that the packages can be set to the destination node when the source lack the information about the destination node’s grouping ID. AODV is a table-driven routing protocol, which will set up the routing table first and then search the effective relay node in the table to transmit the package. The routing protocols above have the relative good performance in MANET and VANET. However, they cannot fit the earthquake emergency scenario. Because these routing protocols do not consider the disaster areas’ distribution character and mobility model’s influence on the rescue quality.\n\n## Conclusion\n\nRouting protocols for post-earthquake emergency communication network are different with traditional ones because of the heterogeneity and dynamicity of the network. After the earthquake, rescue urgency degree is related to the affected areas. In this paper, we first divide the whole disaster area into several regions with different RUD values according to catastrophic intensity obtained. A four-quadrant mobility model (FQMM) for rescuers based on RUD is proposed. With FQMM, the most seriously affected areas are rescued as soon as possible, furthermore, the rescue crews are allocated to each appropriate quadrant, which improves the rescue efficiency and shortens the time of rescue. Under this mobility model, we propose the FQMMBP protocol for emergency communication network, which improves the RREQ message by adding four new fields: RUD, QL, QID and QW.\n\nA virtual earthquake scenario is simulated, and results show that FQMMBP is superior to traditional routing protocols (AODV, DSDV and DSR) in performances of PDR and Delay. We trade 46.88% overhead loss to gain the benefits of 64.45% delay and 16.31% FDR. It is a remarkable and acceptable results for the specific earthquake scenario, which the victims’ life is much more important than rescuers’ battery.\n\nIn our further work, the model will focus on the global rescue, which means that one task may be finished by several rescuers. Once one rescuer exhausts out, the nearby rescuers can receive this information based on the ECNs in this paper and continues the remaining rescue task, which depends on the optimization model.\n\n## Availability of data and materials\n\nData sharing not applicable to this article as no datasets were generated or analyzed during the current study.\n\n## Abbreviations\n\nECV:\n\nThe emergency communication vehicle\n\nBS:\n\nBase station\n\nUAV:\n\nUnnamed aerial vehicle\n\nRUD:\n\nRescue urgency degree, which is inversely proportional to the seismic intensity value in this area. 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Eng. Technol. 4(4), 242–248 (2017)\n\n20. B. Ramakrishnan, R.B. Nishanth, M.M. Joe, M. Selvi, Cluster based emergency message broadcasting technique for vehicular ad hoc network. Wirel. Netw. 23(1), 233–248 (2017). https://doi.org/10.1007/s11276-015-1134-6\n\n21. T. Camp, J. Boleng, V. Davies, A survey of mobility models for ad hoc network research. Wirel. Commun. Mob. Comput. 2(5), 483–502 (2002). https://doi.org/10.1002/wcm.72\n\n22. X. Wang, D. Li, X. Zhang, Y. Cao, MCDM-ECP: multi criteria decision making method for emergency communication protocol in disaster area wireless network. Appl. Sci. 8(7), 1165 (2018). https://doi.org/10.3390/app8071165\n\n23. X. Wang, D. Li, C. Guo, X. Zhang, S.S. Kanhere, K. Li, E. Tovar, Eavesdropping and jamming selection policy for suspicious UAVs based on low power consumption over fading channels. Sensors 19(5), 1126 (2019)\n\n## Acknowledgements\n\nWe would like to thank Ping Li who guided and gave us constructive suggestions at the beginning of the research. This research was funded by China Earthquake Administration under Grant No.XH21009. The preparation of this paper has also received support from Shanghai Sheshan National Geophysical Observatory under Grant No.2020Z04, Technology Commission of Shanghai Municipality under Grant No. 18DZ1200500, and Shanghai Sailing Program under Grant No.21YF1432800.\n\n### Authors' information\n\nXiaoming Wang: received his M.S. degree from Lanzhou University, Lanzhou, China, in 2008. He is currently working as a senior engineer in Shanghai Earthquake Administration; meanwhile, he is working towards his Ph.D. degree in Intelligence Communication in Donghua University, Shanghai, China His research interest includes vehicular networks, sensor networks and emergency communications.. Chang Guo: received the B.S. degree in Communication engineering from Donghua University in 2014 and Ph.D. in Information and intelligent communication system from Donghua University in 2020. Her research interests include the data delivery delay in VANETs, the traffic information acquisition, short-term traffic flow prediction, real-time path planning and load balance optimization in urban area.\n\n## Funding\n\nThis research was funded by China Earthquake Administration under Grant No.XH21009. The preparation of this paper has also received support from Shanghai Sheshan National Geophysical Observatory under Grant No.2020Z04, Technology Commission of Shanghai Municipality under Grant No. 18DZ1200500 and Shanghai Sailing Program under Grant No.21YF1432800.\n\n## Author information\n\nAuthors\n\n### Contributions\n\nXW contributed on the work design, acquisition and analysis of data, and drafted the manuscript. CG did the works on the design of model and algorithms, and revision the manuscript based on the comments. All authors read and approved the final manuscript.\n\n### Corresponding author\n\nCorrespondence to Chang Guo.\n\n## Ethics declarations\n\n### Competing interest\n\nThe project is partially funded by China Earthquake Administration, other parts are paid directly by projects partners. One of the main aims of the activities is the generation of UAV’s FQBMMBP routing protocol. So, it is given to publish all the main funding in the project to force standardization in this area.", null, "" ]
[ null, "https://jwcn-eurasipjournals.springeropen.com/track/article/10.1186/s13638-021-02017-y", null ]
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https://www.gradesaver.com/textbooks/math/algebra/intermediate-algebra-for-college-students-7th-edition/chapter-5-section-5-5-factoring-special-forms-exercise-set-page-371/34
[ "## Intermediate Algebra for College Students (7th Edition)\n\n$2y(1-x^3y)(1+x^3y)$\nRECALL: A difference of two squares can be factored using the formula: $m^2-n^2=(m-n)(m+n)$ Factor out the greatest common factor, $2y$, to obtain: $=2y(1-x^6y^2) \\\\=2y[1^2-(x^3y)^2]$ The binomial is a difference of two squares. Factor the given difference of two squares using the formula above with $m=1$ and $n=x^3y$ to obtain: $=2y(1-x^3y)(1+x^3y)$" ]
[ null ]
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https://www.premiumessays.net/sample-paper-on-difference-between-p-value-of-one-tailed-and-two-tailed-test/
[ "# Sample Paper on Difference between P-Value of One-tailed and Two-tailed Test\n\nWhen conducting a statistical significance test, p-values are generally provided in the output of the population selected. Therefore, ANOVA, regression, and correlation majorly apply the formula during hypothesis testing to assist in understanding the data transmitted. A one-tailed test is one-sided on the critical distribution area, and it is greater than or less than a particular figure, but not both (Hick 316). However, the two-tailed tests is a technique where the critical region is two-sided and confirms if the samples are within the range of values identified in a null hypothesis.\n\nFirst, alpha levels assist in providing relatively accurate results. Hence, it is the probability of arriving an incorrect decision when the null hypothesis is correct. Therefore, one-tailed uses the aggregate value, and it can only appear on either the left or right side. However, a two-tailed splits the alpha into half.  The distinction between and two-tailed P-value is more straightforward in context. In the interpretation of a one-tail P value, it is possible to interpret the group that will have extensive data (Mourougan Sendi and Sethurama 37). Besides, a two-tailed selection involves a random range of benefits from the sample in case the null hypothesis is correct.\n\nIn differentiating the two tests, power must also be considered because it provides a relevant distinction. One-tailed tests subject the researcher into focusing on one side of the distribution table. Moreover, there is more power to one-tiled tests because people are always sure of the decisions made. The two-tailed test considers the probable outcomes making it difficult to rely on one provision. Therefore, it is easier to make conclusions of the sample with the one-tail null hypothesis being correct.\n\nWork Cited\n\nHick, W. E. “A note on one-tailed and two-tailed tests.” (1952): 316.\n\nMourougan, Sendil, and K. Sethuraman. “Hypothesis Development and Testing.” J. Bus. Manag 19 (2017): 34-40." ]
[ null ]
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https://codereview.stackexchange.com/questions/57943/determining-maximum-profit-to-be-made-from-selling-shares
[ "# Determining maximum profit to be made from selling shares\n\nQuestion\n\nYour algorithms have become so good at predicting the market that you now know what the share price of Wooden Orange Toothpicks Inc. (WOT) will be for the next $N$ days.\n\nEach day, you can either buy one share of WOT, sell any number of shares of WOT that you own, or not make any transaction at all. What is the maximum profit you can obtain with an optimum trading strategy?\n\nInput\n\nThe first line contains the number of test cases $T$. $T$ test cases follow:\n\nThe first line of each test case contains a number $N$. The next line contains $N$ integers, denoting the predicted price of WOT shares for the next $N$ days.\n\nOutput\n\nOutput $T$ lines, containing the maximum profit which can be obtained for the corresponding test case.\n\nConstraints\n\n$1 <= T <= 10$\n$1 <= N <= 50000$\n\nAll share prices are between 1 and 100000\n\nCode\n\nimport java.io.*;\nimport java.util.*;\nimport java.text.*;\nimport java.math.*;\nimport java.util.regex.*;\nclass Solution{\npublic static void main(String[] args){\nScanner stdin = new Scanner(System.in);\nint t = stdin.nextInt();\nfor(int m=0;m<t;m++){\nlong n = stdin.nextLong();\nlong[] ar = new long[(int)n];\nfor(int i=0;i<n;i++){\nar[i] = stdin.nextLong();\n}\nlong max = maximum(ar);\nSystem.out.println(max);\n}\n}\npublic static long maximum(long[] ar){\nint j=0;\nlong costPrice=0;\nlong sellingPrice=0;\nint k=0;\nlong max = maximumKey(j,ar);\nwhile(j<ar.length){\nif(ar[j]<max){\ncostPrice += ar[j];\nj++;\n}\nelse if(ar[j]==max){\nsellingPrice += ((j-k)*ar[j]);\nj++;\nk=j;\nmax = maximumKey(j,ar);\n}\n}\nreturn sellingPrice-costPrice;\n}\npublic static long maximumKey(int j, long[] ar){\nlong max = 0;\nfor(int i=j;i<ar.length;i++){\nif(max<ar[i])\nmax = ar[i];\n}\nreturn max;\n}\n}\n\n\nIn the above code I didn't use DP. But this is supposed to be solved using DP with time complexity $O(n)$.\n\n• There is a well known way to do this O(N) time with dynamic if you could only hold on to 1 share at a time. However, I do not think that it is possible in this case (with as many shares as can be purchased), since the most profit would stem from buying all the way until the last global maximum (if it is duplicated), waiting until prices have fallen to the level of the next peak, and recursing. Repeatedly finding the global maximum will cause this to be a O(N^2) problem in the worst case (with constantly rising and falling prices, and diminishing maximums). Jul 24 '14 at 23:57\n• Yes my program is of n^2 level whearas as i saw people saying they were able to solve in O(n) time by using DP. Now how can this question be solved using DP hackerrank.com/challenges/stockmax Jul 25 '14 at 1:15\n• Ah, there is indeed a special trick for determining next global max in O(1), see my answer below. Jul 25 '14 at 5:06\n\nThe trick to getting this problem in O(n) is some clever pre-processing.\n\n# Examples\n\nLet's take 2 trials:\n\nTrial 1: 1, 2, 3, 4, 5\n\nTrail 2: 5, 4, 3, 2, 1\n\nObviously best case in Trial 1 is to buy for 4 days and sell on the 5th, for profit of 10. In Trial 2, there is no profit to be gained, because the price never increases, because there is never a higher maximum down the line.\n\n# Trick\n\nThis points towards the trick: starting from the end make a note of the maximum encountered thus far.\n\nFor Trial 1, the maximums are 5, 5, 5, 5, 5\n\nFor Trial 2, the maximums are the same as the prices\n\nLet's try this for a more complex scenario: 10, 1, 9, 2, 8, 3, 7\n\nThe corresponding maximums would be: 10, 9, 9, 8, 8, 7, 7\n\nThen, for every day that price < maximum, you should buy. If price = maximum, you should sell, and if price > maximum, wait it out.\n\n# Code\n\nThis yields the following code, assuming you parse the prices into a non-empty long[] named prices:\n\nlong[] maximums = new long[prices.length];\nmaximums[prices.length - 1] = prices[prices.length - 1];\n// fill maximums array, from the end\nfor (int i = prices.length - 2; i >= 0; i--) {\nif (prices[i] > maximums[i + 1])\nmaximums[i] = prices[i];\nelse\nmaximums[i] = maximums[i + 1];\n}\n\nlong profit = 0;\nfor (int i = 0; i < prices.length; i++) {\nif (prices[i] < maximums[i]) {\n// BUY now and SELL at max\nprofit += maximums[i] - prices[i];\n}\n}\n\n• Got it.I was calculating the maximum everytime inside the for loop,instead i should have made an array storing the maximum value. Jul 25 '14 at 16:36\n• @mleyfman Thanks your solution helped me.But i guess there is a slight mistake.if(price[i]>maximums[i-1]) should be if price[i]>maximums[i+1],and maximums[i]=maximums[i-1] should be maximums[i]=maximums[i+1]. Correct me if i am wrong Jun 19 '15 at 7:43\n• Thanks for the trick and code. The idea is correct, however, some small typos are in your code which needed to be corrected when you fill out the maximums from the end: if (prices[i] > maximums[i + 1]) maximums[i] = prices[i]; else maximums[i] = maximums[i + 1];\n– user103564\nApr 22 '16 at 6:50\n• @naz20z Thanks for pointing that out, updated answer to reflect it. Apr 30 '16 at 22:38\n• This can be optimized to be single-pass constant space: long max = Long.MIN_VALUE, profit = 0; for(int i = prices.length - 1; i >= 0; --i) if(prices[i] >= max) max = prices[i]; else profit += max - prices[i]; Feb 14 '17 at 22:30" ]
[ null ]
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https://proofwiki.org/wiki/Sum_of_Sines_of_Multiples_of_Angle
[ "# Sum of Sines of Multiples of Angle\n\n## Theorem\n\n $\\displaystyle \\sum_{k \\mathop = 1}^n \\sin k x$ $=$ $\\displaystyle \\sin x + \\sin 2 x + \\sin 3 x + \\cdots + \\sin n x$ $\\displaystyle$ $=$ $\\displaystyle \\frac {\\sin \\frac {\\paren {n + 1} x} 2 \\sin \\frac {n x} 2} {\\sin \\frac x 2}$\n\nwhere $x$ is not an integer multiple of $2 \\pi$.\n\n## Proof 1\n\n$2 \\sin \\alpha \\sin \\beta = \\map \\cos {\\alpha - \\beta} - \\map \\cos {\\alpha + \\beta}$\n\nThus we establish the following sequence of identities:\n\n $\\displaystyle 2 \\sin x \\sin \\frac x 2$ $=$ $\\displaystyle \\cos \\frac x 2 - \\cos \\frac {3 x} 2$ $\\displaystyle 2 \\sin 2 x \\sin \\frac x 2$ $=$ $\\displaystyle \\cos \\frac {3 x} 2 - \\cos \\frac {5 x} 2$ $\\displaystyle$ $\\cdots$ $\\displaystyle$ $\\displaystyle 2 \\sin n x \\sin \\frac x 2$ $=$ $\\displaystyle \\cos \\frac {\\paren {2 n - 1} x} 2 - \\cos \\frac {\\paren {2 n + 1} x} 2$\n\nSumming the above:\n\n $\\displaystyle 2 \\sin \\frac x 2 \\paren {\\sum_{k \\mathop = 1}^n \\sin k x}$ $=$ $\\displaystyle \\cos \\frac x 2 - \\cos \\frac {\\paren {2 n + 1} x} 2$ Sums on right hand side form Telescoping Series $\\displaystyle$ $=$ $\\displaystyle -2 \\map \\sin {\\dfrac {\\frac x 2 + \\frac {\\paren {2 n + 1} x} 2} 2} \\map \\sin {\\dfrac {\\frac x 2 - \\frac {\\paren {2 n + 1} x} 2} 2}$ Prosthaphaeresis Formula for Cosine minus Cosine $\\displaystyle$ $=$ $\\displaystyle -2 \\sin \\dfrac {\\paren {n + 1} x} 2 \\sin \\dfrac {-n x} 2$ $\\displaystyle$ $=$ $\\displaystyle 2 \\sin \\dfrac {\\paren {n + 1} x} 2 \\sin \\dfrac {n x} 2$ Sine Function is Odd\n\nThe result follows by dividing both sides by $2 \\sin \\dfrac x 2$.\n\nIt is noted that when $x$ is a multiple of $2 \\pi$ then:\n\n$\\sin \\dfrac x 2 = 0$\n\nleaving the right hand side undefined.\n\n$\\blacksquare$\n\n## Proof 2\n\nLet $x$ be a real number that is not a integer multiple of $2 \\pi$.\n\nLet $k$ be a non-negative integer.\n\nWe have, from Euler's Formula:\n\n$\\map \\exp {i k x} = i \\sin k x + \\cos k x$\n\nSumming from $k = 0$ to $k = n$, we have:\n\n$\\displaystyle \\sum_{k \\mathop = 0}^n \\map \\exp {i k x} = i \\sum_{k \\mathop = 0}^n \\sin k x + \\sum_{k \\mathop = 0}^n \\cos k x$\n\nAs $\\sin k x$ and $\\cos k x$ are both real for real $k, x$, we have:\n\n $\\displaystyle \\sum_{k \\mathop = 0}^n \\sin k x$ $=$ $\\displaystyle \\map \\Im {\\sum_{k \\mathop = 0}^n \\map \\exp {i k x} }$ $\\displaystyle$ $=$ $\\displaystyle \\map \\Im {\\paren {i \\sin \\frac {n x} 2 + \\cos \\frac {n x} 2} \\frac {\\map \\sin {\\frac {\\paren {n + 1} x} 2} } {\\sin \\frac x 2} }$ Sum of $\\map \\exp {i k x}$ $\\displaystyle$ $=$ $\\displaystyle \\frac {\\sin \\frac {\\paren {n + 1} x} 2 \\sin \\frac {n x} 2} {\\sin \\frac x 2}$\n\n$\\blacksquare$" ]
[ null ]
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https://codecrucks.com/bubble-sort/
[ "# Bubble Sort\n\nBubble sort is comparison based sorting method, and also known as sinking sort. It is perhaps most simple sorting algorithm.\n\nBubble Sort is a simple method for sorting a given set of n elements provided in the form of an array with n elements. It analyzes each element individually and sorts them based on their values.\n\nIt compares the first and second elements of the array; if the first element is greater than the second element, it will swap both elements, and then compare the second and third elements, and so on.\n\nIn general, sorting is accomplished by comparing neighboring elements A[j] and A[j+1]. If the components being compared are out of sequence, they are switched. The biggest element bubbles up at the final place in the unsorted array at the end of each iteration. For array of size n, this process is repeated n – 1 times.\n\nSorting an array starts from the end. The largest element is moved to the last place in the first iteration, the second largest element is moved to the second last position in the second iteration, and so on.\n\nIt is called as bubble sort because, with each complete iteration, the largest element in the given array bubbles up to the last position or the highest index, similar to how a water bubble rises to the water’s surface.\n\nThis basic method performs badly in practice and is mostly used as an educational tool. Sorting libraries integrated into popular computer languages like as Python and Java utilize more efficient algorithms such as quicksort, timsort, or merge sort.\n\n## Example of Bubble Sort\n\nBubble sort compares adjacent elements and swaps them if they are out of order. In every iteration, the largest element from unsorted array moves to the last location. Following figures show the step by step simulation of it.\n\nlet us sort the letters of word “DESIGN” in alphabetical order using bubble sort.\n\nStep by step comparison and output of every pass is depicted in following figures. Adjacent cells in light gray colors are under comparison and dark gray cells indicate sorted elements.\n\n## Algorithm for Bubble Sort\n\nAlgorithm BUBBLE_SORT(A)\n// A is an array of size n\n\nfor i ← 1 to n do\nfor j ← 1 to n – i do\nif A[j] > A[j+1] do\nswap(A[j], A[j+1])\nend\nend\nend\n\nAlthough the above logic would sort an unsorted array, the technique is inefficient since the outer for loop will continue to execute for n iterations even if the array is sorted.\n\nOptimized version of the algorithm is mentioned below:\n\nAlgorithm OPTIMIZED_BUBBLE_SORT(A)\n// A is an array of size n\n\nfor i ← 1 to n do\nflag ← 1\nfor j ← 1 to n – i do\nif A[j] > A[j+1] do\ntemp ← A[j]\nA[i] ← A[j+1]\nA[j+1] ← temp\nflag ← 0\nend\nif (flag) do\nbreak\nend\nend\nend\n\n## Complexity Analysis\n\nFor first iteration of outer loop, inner loop does n comparisons. For second iteration of outer loop, inner loop dose n – 1 comparison and so on. In last iteration of outer loop, inner loop does only one comparison. Outer loop itself iterates n times. Running time of algorithm is defined as total number of comparisons. Thus,\n\nDespite the fact that statements within the inner loop do not execute, the complexity of bubble sort is O(n2). The quadratic complexity is caused by the two nested for loops, which iterate regardless of the input data pattern. As a result, the complexity of bubble sort is the same whether the situation is best, worst, or average.\n\nIt’s self-explanatory, that if no swapping occurs in the first run, the list is already sorted. By setting an appropriate flag, we can break the cycle. In this scenario, the best-case bubble sort running time would be linear, i.e. O(n).\n\nDuring the first iteration of the outer loop, the inner loop executes n times without changing the flag, indicating that the data is sorted.\n\nFollowing table shows the time complexity of bubble sort for both versions, normal (without flag) and optimized (with flag):\n\nOther O(n2) sorting algorithms, such as insertion sort, are typically quicker than bubble sort and are not any more complex. As a result, bubble sort is not a useful sorting algorithm.\n\nThe single important benefit it has over most other algorithms, including quicksort but not insertion sort, is that the ability to discover if list is sorted. The complexity of this algorithm is only O(n) when the list is already sorted (best-case).\n\nMost other algorithms, even ones with lower average-case complexity, execute their whole sorting operation on the set, making them more complicated.\n\n## Simulation\n\nWorking of bubble sort is nicely simulated in following diagram\n\n## Rabbits and Turtles\n\nBecause elements flow in different directions at different speeds, the distance and direction that elements must move during the sort influence the performance of bubble sort.\n\nLarger element may participate in consecutive swaps and moves faster toward end. smaller elements moves quite slowly towards beginning. Consider following data patter, where the largest element is at beginning and the smallest is at end.\n\nA = <8, 7, 6, 5, 4, 3, 2, 1>\n\nAfter first iteration, array A would be,\n\nA = <7, 6, 5, 4, 3, 2, 1, 8>\n\nAs can be seen, in one iteration, 8 moved from first to last position (moved seven positions.) Whereas, element 1 moved from last to second last position (moved one position only).\n\nSo larger element, which moved quickly towards the end are known as hare, and smaller elements which moves slowly towards beginning are known as turtle.\n\n## Discussion and Comparison\n\nSeveral attempts have been made to eliminate turtles in order to increase the speed of bubble sorting. Cocktail sort is a bidirectional bubble sort that moves from beginning to end and then reverses direction, moving from end to beginning. It can move turtles rather effectively, but its worst-case complexity is O(n2).\n\nIn his famous book “The Art of Computer Programming”, Donald Knuth concludes that “the bubble sort appears to have little to recommend it, except a catchy name and the fact that it leads to some interesting theoretical issues\n\nComb sort compares elements separated by large gaps and can move turtles fast before moving on to smaller and smaller gaps to smooth down the list. Its average performance is comparable to that of quicker algorithms such as quicksort.\n\nIn the worst situation, bubble sort is asymptotically identical to insertion sort in terms of running time, but the two algorithms differ significantly in terms of the number of swaps required. Astrachan’s experimental results also demonstrate that insertion sort works significantly better even on random lists. For these reasons, many current algorithm textbooks prefer the insertion sort method over the bubble sort technique.\n\nBubble sort also has a bad interaction with current CPU technology. It generates at least twice the number of writes as insertion sort, double the number of cache misses, and asymptotically more branch mispredictions. Astrachan’s experiments sorting strings in Java demonstrate that bubble sort is about one-fifth as quick as an insertion sort and 70% faster than a selection sort.\n\n## Code for Bubble Sort\n\nFollowing are the few popular sorting algorithms used in computer science:\n\n### 2 Responses\n\n1.", null, "paresh says:\n\n•", null, "codecrucks says:" ]
[ null, "https://secure.gravatar.com/avatar/6be4caaf700cc7a2b66503633b54a0ca", null, "https://secure.gravatar.com/avatar/390e4c6d856849f1ace10209aca7584b", null ]
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https://www.isrgrajan.com/equations-you-should-know-that-changed-the-future-of-physics.html
[ "Thursday, September 28, 2023\nEducationEquations You Should Know: That Changed The Future Of Physics\n\n# Equations You Should Know: That Changed The Future Of Physics\n\nIt is well known that science without physics is nothing, but lame. Even physics can be considered the backbone of science. Every great invention is interlinked to physics. Physics can be divided majorly in two sections that is classical and modern physics. However, classical is very old one and deals with mainly person’s like Galileo and Archimedes, but however we are going to sketch only modern physics and the theories that developed and changed the whole physics. There are many laws dealing the same, but we are going to know about the basic ones which is of importance in the view of instruments and devices which we use in our daily life or the theories which developed and helped in broadening our minds thinking about outside of the earth.\n\n### 1. Newton’s Law Of Universal Gravitation:\n\nThe law and theory were developed by the great scientist Issac Newton. It got published in Printipia in July 1687 for the first time. The law is very special because it deals with the facts of revolving planet’s around the sun in the orbit and the factors affecting the revolution. Secondly, the law deals with the gravity of the revolving planets. This is the main section which is the basis for the Aerospace education and helps in mainly finding out all the concepts required to eject the satellites outside the space.\n\n### 2. Theory Of Relativity:\n\nThis law was proposed by Albert Einstein in 1905. This is the basic equation taught to every high school student of physics, the equation which is taught to students is E=mc^2. Why it is so important? This is the most accepted theory on the relationship between space and earth, it was derived after around two hundred years of Newton’s law of gravitation. Most of the Knowledge which we have about space and the universe is known with this Einstein’s great relation. Well, after he proposed this theory many other theories got developed with respect to relativity accordingly, but thus the basic equation which helps everyone who is presenting developed one.\n\n### 3. Maxwell’s Equations:\n\nThis law was proposed by James Clark Maxwell. Maxwell proposed this law in around 1860’s. This law deals with how charged particles give rise to electric and magnetic field and by this time you would have understood why it is so important to study this topic? This is basic theory which is helpful in running the great trains, running of different motors and magnetic instruments. Maxwell’s equations are the primary syllabus of the electrical and electronics graduates they basically have to study this equation in their entire curriculum. The basic application is observed in our induction cooking, ever wondered on which principle it runs, obviously it is a Maxwells law of electromagnetic induction which is monitoring it.\n\n### 4. Law’s Of Thermodynamics:\n\nThere are three law’s of Thermodynamics but second one plays a key role. This law was developed by Rudolf Clausius. This law deals with how energy flows from higher concentration to lower concentration, it was developed in 1865. This law is responsible for the development of the instruments and appliances which are used in industries, household and everywhere we look at. For example the refrigerator Which we use in our house is developed according to this law, motors, compressors everything runs on this principle. When dealing with studies this law is useful for the students of mechanical engineering.\n\n### 5. Schrödinger’s Equation:\n\nThis equation was proposed by Austrian physicist Erwin Schrödinger in 1926. This law is the basic equation of the area quantum physics. This law, mainly deals with how the quantum state of a quantum system with respect to time. The law helped in developing various instruments including electron microscope, microchips and lot more. This law helped in widening the facts and myths related to atoms and it’s constituents present inside it.\n\n### 6. Newton’s law of motion:\n\nEvery one of us are aware about the incident took place with Newton which led to the discovery of gravity. Yes, you are right falling of an apple. Most of them include this topic in classical physics but it is almost helpful in modern physics equally. These laws are about three in number, first one speaks about inertia, the second deals with momentum and the last one deals with the reaction force. The laws are fundamental for all the Engineers.\n\n### 7. Logaritms:\n\nWhen it comes to logarithms it is the most accepted practice when a bigger calculation has to be done. It was developed by Scottish Laird, John Napier of Merchiston. When it comes to solving multiplication of large numbers, for example 3456×56789 is really time consuming and needs good skills. To overcome this problem Napier came up with logarithms which us boon for calculations. Now, due to the advent of computers and calculators is of less use because it is inbuilt whereas, scientists always use logarithms in their day to day life. The equations are mainly used in finding exponential equations, compound interest and radioactive decay by a Scientist, Astronimer and an Engineer.\n\n### 8. Coloumb’s Law Of Electrostatics:\n\nWhen it comes to electricity electrostatics is the main branch dealing it and in turn coulombs law is the basic law of electrostatics. This law explains how force is related to charges and the distance between them. Interestingly, it helped in developing many electrical instruments which we use nowadays from fan to its capacitors, generators and what not.\n\nLast but not least there are many laws are proposed by scientist’s day by day, each law has got its own importance, it’s own dominancy, its own applications and every law is of great importance. But the major laws which led to great transformation in our society are the one discussed above. If you are a science student it is obvious that you may know all these laws, but it is exactly not like that even being general everyone should know about this law at least basic terminology because every appliance which we use has a principle due to which it is working and as we are using it we have to know the basic law which is governing that appliance or instrument." ]
[ null ]
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http://docs.juliadiffeq.org/latest/tutorials/ode_example.html
[ "Ordinary Differential Equations\n\n# Ordinary Differential Equations\n\nThis tutorial will introduce you to the functionality for solving ODEs. Other introductions can be found by checking out DiffEqTutorials.jl. Additionally, a video tutorial walks through this material.\n\n## Example 1 : Solving Scalar Equations\n\nIn this example we will solve the equation\n\n$\\frac{du}{dt} = f(u,p,t)$\n\non the time interval $t\\in[0,1]$ where $f(u,p,t)=αu$. We know by Calculus that the solution to this equation is $u(t)=u₀\\exp(αt)$.\n\nThe general workflow is to define a problem, solve the problem, and then analyze the solution. The full code for solving this problem is:\n\nusing DifferentialEquations\nf(u,p,t) = 1.01*u\nu0=1/2\ntspan = (0.0,1.0)\nprob = ODEProblem(f,u0,tspan)\nsol = solve(prob,Tsit5(),reltol=1e-8,abstol=1e-8)\nusing Plots\nplot(sol,linewidth=5,title=\"Solution to the linear ODE with a thick line\",\nxaxis=\"Time (t)\",yaxis=\"u(t) (in μm)\",label=\"My Thick Line!\") # legend=false\nplot!(sol.t, t->0.5*exp(1.01t),lw=3,ls=:dash,label=\"True Solution!\")\n\nwhere the pieces are described below.\n\n### Step 1: Defining a Problem\n\nTo solve this numerically, we define a problem type by giving it the equation, the initial condition, and the timespan to solve over:\n\nusing DifferentialEquations\nf(u,p,t) = 1.01*u\nu0=1/2\ntspan = (0.0,1.0)\nprob = ODEProblem(f,u0,tspan)\n\nNote that DifferentialEquations.jl will choose the types for the problem based on the types used to define the problem type. For our example, notice that u0 is a Float64, and therefore this will solve with the dependent variables being Float64. Since tspan = (0.0,1.0) is a tuple of Float64's, the independent variables will be solved using Float64's (note that the start time and end time must match types). You can use this to choose to solve with arbitrary precision numbers, unitful numbers, etc. Please see the notebook tutorials for more examples.\n\nThe problem types include many other features, including the ability to define mass matrices and hold callbacks for events. Each problem type has a page which details its constructor and the available fields. For ODEs, the appropriate page is here. In addition, a user can specify additional functions to be associated with the function in order to speed up the solvers. These are detailed at the performance overloads page.\n\n### Step 2: Solving a Problem\n\n#### Controlling the Solvers\n\nAfter defining a problem, you solve it using solve.\n\nsol = solve(prob)\n\nThe solvers can be controlled using the available options are described on the Common Solver Options manual page. For example, we can lower the relative tolerance (in order to get a more correct result, at the cost of more timesteps) by using the command reltol:\n\nsol = solve(prob,reltol=1e-6)\n\nThere are many controls for handling outputs. For example, we can choose to have the solver save every 0.1 time points by setting saveat=0.1. Chaining this with the tolerance choice looks like:\n\nsol = solve(prob,reltol=1e-6,saveat=0.1)\n\nMore generally, saveat can be any collection of time points to save at. Note that this uses interpolations to keep the timestep unconstrained to speed up the solution. In addition, if we only care about the endpoint, we can turn off intermediate saving in general:\n\nsol = solve(prob,reltol=1e-6,save_everystep=false)\n\nwhich will only save the final time point.\n\n#### Choosing a Solver Algorithm\n\nDifferentialEquations.jl has a method for choosing the default solver algorithm which will find an efficient method to solve your problem. To help users receive the right algorithm, DifferentialEquations.jl offers a method for choosing algorithms through hints. This default chooser utilizes the precisions of the number types and the keyword arguments (such as the tolerances) to select an algorithm. Additionally one can provide alg_hints to help choose good defaults using properties of the problem and necessary features for the solution. For example, if we have a stiff problem where we need high accuracy, but don't know the best stiff algorithm for this problem, we can use:\n\nsol = solve(prob,alg_hints=[:stiff],reltol=1e-8,abstol=1e-8)\n\nYou can also explicitly choose the algorithm to use. DifferentialEquations.jl offers a much wider variety of solver algorithms than traditional differential equations libraries. Many of these algorithms are from recent research and have been shown to be more efficient than the \"standard\" algorithms. For example, we can choose a 5th order Tsitouras method:\n\nsol = solve(prob,Tsit5())\n\nNote that the solver controls can be combined with the algorithm choice. Thus we can for example solve the problem using Tsit5() with a lower tolerance via:\n\nsol = solve(prob,Tsit5(),reltol=1e-8,abstol=1e-8)\n\nIn DifferentialEquations.jl, some good \"go-to\" choices for ODEs are:\n\n• AutoTsit5(Rosenbrock23()) handles both stiff and non-stiff equations. This is a good algorithm to use if you know nothing about the equation.\n• BS3() for fast low accuracy non-stiff.\n• Tsit5() for standard non-stiff. This is the first algorithm to try in most cases.\n• Vern7() for high accuracy non-stiff.\n• Rodas4() for stiff equations with Julia-defined types, events, etc.\n• radau() for really high accuracy stiff equations (requires installing ODEInterfaceDiffEq.jl)\n\nFor a comprehensive list of the available algorithms and detailed recommendations, Please see the solver documentation. Every problem type has an associated page detailing all of the solvers associated with the problem.\n\n### Step 3: Analyzing the Solution\n\n#### Handling the Solution Type\n\nThe result of solve is a solution object. We can access the 5th value of the solution with:\n\nsol #.637\n\nor get the time of the 8th timestep by:\n\nsol.t #.438\n\nConvenience features are also included. We can build an array using a comprehension over the solution tuples via:\n\n[t+u for (u,t) in tuples(sol)]\n\nor more generally\n\n[t+2u for (u,t) in zip(sol.u,sol.t)]\n\nallows one to use more parts of the solution type. The object that is returned by default acts as a continuous solution via an interpolation. We can access the interpolated values by treating sol as a function, for example:\n\nsol(0.45) # The value of the solution at t=0.45\n\nNote the difference between these: indexing with [i] is the value at the ith step, while (t) is an interpolation at time t!\n\nIf in the solver dense=true (this is the default unless saveat is used), then this interpolation is a high order interpolation and thus usually matches the error of the solution time points. The interpolations associated with each solver is detailed at the solver algorithm page. If dense=false (unless specifically set, this only occurs when save_everystep=false or saveat is used) then this defaults to giving a linear interpolation.\n\nFor details on more handling the output, see the solution handling page.\n\n#### Plotting Solutions\n\nWhile one can directly plot solution time points using the tools given above, convenience commands are defined by recipes for Plots.jl. To plot the solution object, simply call plot:\n\n#]add Plots # You need to install Plots.jl before your first time using it!\nusing Plots\n#plotly() # You can optionally choose a plotting backend\nplot(sol)", null, "If you are in Juno, this will plot to the plot pane. To open an interactive GUI (dependent on the backend), use the gui command:\n\ngui()\n\nThe plot function can be formatted using the attributes available in Plots.jl. Additional DiffEq-specific controls are documented at the plotting page.\n\nFor example, from the Plots.jl attribute page we see that the line width can be set via the argument linewidth. Additionally, a title can be set with title. Thus we add these to our plot command to get the correct output, fix up some axis labels, and change the legend (note we can disable the legend with legend=false) to get a nice looking plot:\n\nplot(sol,linewidth=5,title=\"Solution to the linear ODE with a thick line\",\nxaxis=\"Time (t)\",yaxis=\"u(t) (in μm)\",label=\"My Thick Line!\") # legend=false\n\nWe can then add to the plot using the plot! command:\n\nplot!(sol.t,t->0.5*exp(1.01t),lw=3,ls=:dash,label=\"True Solution!\")", null, "## Example 2: Solving Systems of Equations\n\nIn this example we will solve the Lorenz equations:\n\n\\begin{align} \\frac{dx}{dt} &= σ(y-x) \\\\ \\frac{dy}{dt} &= x(ρ-z) - y \\\\ \\frac{dz}{dt} &= xy - βz \\\\ \\end{align}\n\nDefining your ODE function to be in-place updating can have performance benefits. What this means is that, instead of writing a function which outputs its solution, you write a function which updates a vector that is designated to hold the solution. By doing this, DifferentialEquations.jl's solver packages are able to reduce the amount of array allocations and achieve better performance.\n\nThe way we do this is we simply write the output to the 1st input of the function. For example, our Lorenz equation problem would be defined by the function:\n\nfunction lorenz(du,u,p,t)\ndu = 10.0*(u-u)\ndu = u*(28.0-u) - u\ndu = u*u - (8/3)*u\nend\n\nand then we can use this function in a problem:\n\nu0 = [1.0;0.0;0.0]\ntspan = (0.0,100.0)\nprob = ODEProblem(lorenz,u0,tspan)\nsol = solve(prob)\n\nUsing the plot recipe tools defined on the plotting page, we can choose to do a 3D phase space plot between the different variables:\n\nplot(sol,vars=(1,2,3))", null, "Note that the default plot for multi-dimensional systems is an overlay of each timeseries. We can plot the timeseries of just the second component using the variable choices interface once more:\n\nplot(sol,vars=(0,2))", null, "Note that here \"variable 0\" corresponds to the independent variable (\"time\").\n\n## Defining Parameterized Functions\n\nIn many cases you may want to explicitly have parameters associated with your differential equations. This can be used by things like parameter estimation routines. In this case, you use the p values via the syntax:\n\nfunction parameterized_lorenz(du,u,p,t)\ndu = p*(u-u)\ndu = u*(p-u) - u\ndu = u*u - p*u\nend\n\nand then we add the parameters to the ODEProblem:\n\nu0 = [1.0,0.0,0.0]\ntspan = (0.0,1.0)\np = [10.0,28.0,8/3]\nprob = ODEProblem(parameterized_lorenz,u0,tspan,p)\n\nWe can make our functions look nicer by doing a few tricks. For example:\n\nfunction parameterized_lorenz(du,u,p,t)\nx,y,z = u\nσ,ρ,β = p\ndu = dx = σ*(y-x)\ndu = dy = x*(ρ-z) - y\ndu = dz = x*y - β*z\nend\n\nNote that the type for the parameters p can be anything: you can use arrays, static arrays, named tuples, etc. to enclose your parameters in a way that is sensible for your problem.\n\nAdditionally, there exists a @ode_def macro allows for \"defining your ODE in pseudocode\" and getting a function which is efficient and runnable. To use the macro, you write out your system of equations with the left-hand side being d_ and those variables will be parsed as the dependent variables. The independent variable is t, and the other variables are parameters which you pass at the end. For example, we can write the Lorenz system as:\n\n#]add ParameterizedFunctions\nusing ParameterizedFunctions\ng = @ode_def begin\ndx = σ*(y-x)\ndy = x*(ρ-z) - y\ndz = x*y - β*z\nend σ ρ β\n\nDifferentialEquations.jl will automatically translate this to be exactly the same as f. The result is more legible code with no performance loss. For more information on the macro Domain Specific Language (DSL) and its limitations, please see the parameterized function page The result is that g is a function which you can now use to define the Lorenz problem.\n\nu0 = [1.0;0.0;0.0]\ntspan = (0.0,1.0)\np = [10.0,28.0,8/3]\nprob = ODEProblem(g,u0,tspan,p)\n\nThe macro does \"behind-the-scenes\" symbolic calculations to pre-compute things like the Jacobian, inverse Jacobian, etc. in order to speed up calculations. Thus not only will this lead to legible ODE definitions, but \"unfairly fast\" code! We can turn off some of the calculations by using a more specific macro, like @ode_def_bare. See ParameterizedFunctions.jl for more details.\n\nSince the parameters exist within the function, functions defined in this manner can also be used for sensitivity analysis, parameter estimation routines, and bifurcation plotting. This makes DifferentialEquations.jl a full-stop solution for differential equation analysis which also achieves high performance.\n\n## Example 3: Using Other Types for Systems of Equations\n\nDifferentialEquations.jl can handle many different dependent variable types (generally, anything with a linear index should work!). So instead of solving a vector equation, let's let u be a matrix! To do this, we simply need to have u0 be a matrix, and define f such that it takes in a matrix and outputs a matrix. We can define a matrix of linear ODEs as follows:\n\nA = [1. 0 0 -5\n4 -2 4 -3\n-4 0 0 1\n5 -2 2 3]\nu0 = rand(4,2)\ntspan = (0.0,1.0)\nf(u,p,t) = A*u\nprob = ODEProblem(f,u0,tspan)\n\nHere our ODE is on a 4x2 matrix, and the ODE is the linear system defined by multiplication by A. To solve the ODE, we do the same steps as before.\n\nsol = solve(prob)\nplot(sol)", null, "We can instead use the in-place form by using Julia's in-place matrix multiplication function mul!:\n\nusing LinearAlgebra\nf(du,u,p,t) = mul!(du,A,u)\n\nAdditionally, we can use non-traditional array types as well. For example, StaticArrays.jl offers immutable arrays which are stack-allocated, meaning that their usage does not require any (slow) heap-allocations that arrays normally have. This means that they can be used to solve the same problem as above, with the only change being the type for the initial condition and constants:\n\nusing StaticArrays, DifferentialEquations\nA = @SMatrix [ 1.0 0.0 0.0 -5.0\n4.0 -2.0 4.0 -3.0\n-4.0 0.0 0.0 1.0\n5.0 -2.0 2.0 3.0]\nu0 = @SMatrix rand(4,2)\ntspan = (0.0,1.0)\nf(u,p,t) = A*u\nprob = ODEProblem(f,u0,tspan)\nsol = solve(prob)\nusing Plots; plot(sol)\n\nNote that the analysis tools generalize over to systems of equations as well.\n\nsol\n\nstill returns the solution at the fourth timestep. It also indexes into the array as well. The last value is the timestep, and the beginning values are for the component. This means\n\nsol[5,3]\n\nis the value of the 5th component (by linear indexing) at the 3rd timepoint, or\n\nsol[2,1,:]\n\nis the timeseries for the component which is the 2nd row and 1 column.\n\n## Going Beyond ODEs: How to Use the Documentation\n\nNot everything can be covered in the tutorials. Instead, this tutorial will end by pointing you in the directions for the next steps.\n\n#### Common API for Defining, Solving, and Plotting\n\nOne feature of DifferentialEquations.jl is that this pattern for solving equations is conserved across the different types of differential equations. Every equation has a problem type, a solution type, and the same solution handling (+ plotting) setup. Thus the solver and plotting commands in the Basics section applies to all sorts of equations, like stochastic differential equations and delay differential equations. Each of these different problem types are defined in the Problem Types section of the docs. Every associated solver algorithm is detailed in the Solver Algorithms section, sorted by problem type. The same steps for ODEs can then be used for the analysis of the solution.\n\n#### Additional Features and Analysis Tools\n\nIn many cases, the common workflow only starts with solving the differential equation. Many common setups have built-in solutions in DifferentialEquations.jl. For example, check out the features for:\n\nMany more are defined in the relevant sections of the docs. Please explore the rest of the documentation, including tutorials for getting started with other types of equations. In addition, to get help, please either file an issue at the main repository or come have an informal discussion at our Gitter chatroom." ]
[ null, "http://docs.juliadiffeq.org/latest/assets/ode_tutorial_linear_plot.png", null, "http://docs.juliadiffeq.org/latest/assets/ode_tutorial_thick_linear.png", null, "http://docs.juliadiffeq.org/latest/assets/3d_lorenz.png", null, "http://docs.juliadiffeq.org/latest/assets/lorenz_timeseries.png", null, "http://docs.juliadiffeq.org/latest/assets/multiODEplot.png", null ]
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https://deeplearn.org/arxiv/96589/implicit-neural-solver-for-time-dependent-linear-pdes-with-convergence-guarantee
[ "### Implicit Neural Solver for Time-dependent Linear PDEs with Convergence Guarantee\n\n• 2019-10-08 15:20:01\n• Suprosanna Shit, Abinav Ravi, Ivan Ezhov, Jana Lipkova, Marie Piraud, Bjoern Menze\n• 1\n\n### Abstract\n\nFast and accurate solution of time-dependent partial differential equations(PDEs) is of key interest in many research fields including physics,engineering, and biology. Generally, implicit schemes are preferred over theexplicit ones for better stability and correctness. The existing implicitschemes are usually iterative and employ a general-purpose solver which may besub-optimal for a specific class of PDEs. In this paper, we propose a neuralsolver to learn an optimal iterative scheme for a class of PDEs, in adata-driven fashion. We achieve this goal by modifying an iteration of anexisting semi-implicit solver using a deep neural network. Further, we providetheoretical proof that our approach preserves the correctness and convergenceguarantees provided by the existing iterative-solvers. We also demonstrate thatour model generalizes to a different parameter setting than the one seen duringtraining and achieves faster convergence compared to the semi-implicit schemes.\n\n# Implicit Neural Solver for Time-dependent Linear PDEs with Convergence Guarantee\n\nSuprosanna Shit\nTechnical University Munich\nAbinav Ravi\nTechnical University Munich\nIvan Ezhov\nTechnical University Munich\nJana Lipkova\nTechnical University Munich\nMarie Piraud\nKonica Minolta Laboratory Europe\nBjoern Menze\nTechnical University Munich\n###### Abstract\n\nFast and accurate solution of time-dependent partial differential equations (PDEs) is of key interest in many research fields including physics, engineering, and biology. Generally, implicit schemes are preferred over the explicit ones for better stability and correctness. The existing implicit schemes are usually iterative and employ a general-purpose solver which may be sub-optimal for a specific class of PDEs. In this paper, we propose a neural solver to learn an optimal iterative scheme for a class of PDEs, in a data-driven fashion. We achieve this goal by modifying an iteration of an existing semi-implicit solver using a deep neural network. Further, we provide theoretical proof that our approach preserves the correctness and convergence guarantees provided by the existing iterative-solvers. We also demonstrate that our model generalizes to a different parameter setting than the one seen during training and achieves faster convergence compared to the semi-implicit schemes.\n\nImplicit Neural Solver for Time-dependent Linear PDEs with Convergence Guarantee\n\nSuprosanna Shit Technical University Munich Abinav Ravi Technical University Munich Ivan Ezhov Technical University Munich Jana Lipkova Technical University Munich Marie Piraud Konica Minolta Laboratory Europe Bjoern Menze Technical University Munich\n\n\\@float\n\nnoticebox[b]33rd Conference on Neural Information Processing Systems (NeurIPS 2019), Vancouver, Canada.\n\n## 1 Introduction\n\nTime-dependent partial differential equations (PDEs) are an essential mathematical tool to describe numerous physical processes such as heat transfer, wave propagation, quantum transport, and tumor growth. Solving the initial-value problem (IVP) and boundary-value problem (BVP) accurately and efficiently is of primary research interest in computational science. Numerical solution of time-dependant PDEs requires appropriate spatio-temporal discretization. Spatial discretization can be cast as either finite difference method (FDM) or finite element method (FEM) or finite volume method (FVM), whereas temporal discretization relies on either explicit, implicit or semi-implicit methods. Explicit temporal update rules are generally a single or few forward computation steps, while implicit or semi-implicit update rules, such as Crank-Nicolson’s scheme, resort to a fixed-point iterative scheme. Small time and spatial resolution facilitate a more accurate solution, however, increases computational burden at the same time. Moreover, maximally allowed spatio-temporal resolution is not only constrained by the desired accuracy but also limited to numerical stability criteria. Note that implicit and semi-implicit methods offer a relaxed stability constraints (sometimes unconditionally stable) at the expense of an increased computational cost caused by the iterative schemes.\n\nIn recent times, deep neural networks have gained significant attention by numerical computation community due to its superior performance in solving forward simulations and inverse-problems . Recent work by Tompson et al. shows a data-driven convolutional neural network (CNN) can accelerate fluid simulation compared to traditional numerical schemes. Long et al. shows that learned differential operator can outperform hand-designed discrete schemes. However, on the contrary to the well understood and theoretically grounded classical methods, the deep learning-based approaches rely mainly on empirical validity. Recently, Hsieh et al. develop a promising way to learn numerical solver while providing a theoretical convergence guarantee. They demonstrate that a feed-forward CNN, which is trained to mimic a single iteration of a linear solver, can deliver faster solution than the handcrafted solver. Astonishingly for time-dependent PDEs, the temporal update step of the previously proposed neural schemes relies on an explicit forward Euler method and none of them is capable of making use of the advanced implicit and semi-implicit methods. This limitation restricts the general use of neural architectures in solving time-dependent PDEs.\n\nIn this paper, we introduce a novel neural solver for time-dependant linear PDEs. Motivated by we construct a neural iterator from a semi-implicit update rule. We replace a single iteration of the semi-implicit scheme with a learnable parameterized function such that the fixed point of the algorithm is preserved. To leverage the theoretical guarantees we perform data-driven learning to enhance the convergence speed. As a result, our approach provides: (i) theoretical guarantees of convergence to the correct stationary solution, (ii) faster convergence than existing solvers, and (iii) generalizes to a resolutions and parameter settings very different from the ones seen at training time.\n\n## 2 Methodology\n\nIn this section, we first introduce the necessary background on the semi-implicit method for time-dependant linear PDEs and subsequently describe our proposed neural solver.\n\n### 2.1 Background\n\nIn the following, we consider the general IVP form of time-dependant linear PDE for the variable of interest $u$ within the computation domain $\\Omega$, w.r.t. the spatial variable $x$ and temporal variable $t$, subject to Dirichlet boundary condition $b(x,t)$ at the boundary $\\Gamma$\n\n $\\frac{\\partial u}{\\partial t}=\\mathcal{F}(u,x,t;\\Theta),\\leavevmode\\nobreak\\ % \\forall x\\in\\Omega,\\mbox{ s.t. }u(x,t)=b(x,t),\\leavevmode\\nobreak\\ \\forall x% \\in\\Gamma\\mbox{ and }u_{t_{0}}=u_{0};$ (1)\n\n$\\mathcal{F}(u,x,t;\\Theta)$ is a linear operator and can be discretized as $\\sum_{i=1}^{N}\\frac{\\Theta_{i}\\partial_{i}}{\\delta x^{p_{i}}}u$, with uniform spatial discretization step $\\delta x$ and PDE parameter set $\\Theta=\\{\\Theta_{i}\\}_{i=1:N}$, where $\\Theta_{i}$ is a diagonal matrix comprising the coefficients of the differential operator $\\partial_{i}$ of order $p_{i}$. Let’s call $u(x,t)$ as $u_{t}$ for simplicity. A first order semi-implicit update rule to get $u_{t+\\delta t}$ from $u_{t}$ (with time step $\\delta t$) is given by\n\n $\\frac{u_{t+\\delta t}-u_{t}}{\\delta t}=\\epsilon\\mathcal{F}(u,x,t+\\delta t;% \\Theta)+(1-\\epsilon)\\mathcal{F}(u,x,t;\\Theta);\\leavevmode\\nobreak\\ [0<\\epsilon% \\leq 1]$ (2)\n\nTo obtain $u_{t+\\delta t}$ one needs to solve the following linear system of equations\n\n $\\left(I-\\delta t\\epsilon\\sum_{i=1}^{N}\\frac{\\Theta_{i}d_{i}}{\\delta x^{p_{i}}}% \\right)u_{t+\\delta t}=\\delta t\\epsilon\\sum_{i=1}^{N}\\frac{\\Theta_{i}(\\partial_% {i}-d_{i}I)}{\\delta x^{p_{i}}}u_{t+\\delta t}+c(u_{t},\\Theta,\\delta x,\\delta t,% \\epsilon;\\partial)$ (3)\n\nwhere $c$ is independent of $u_{t+\\delta t}$ and $d_{i}$ is the central element of the central difference discretization of operator $\\partial_{i}$. Note that for central difference scheme, $\\partial_{i}-d_{i}I$ is real, zero-diagonal, and either circulant or skew-circulant matrix. One can use an iterative scheme to compute $u_{t+\\delta t}$ from an arbitrary initialization $u_{t+\\delta t}^{0}$ on the right-hand-side of Eq. 3. For simplicity of notation, we refer to $\\left(I-\\delta t\\epsilon\\sum_{j=1}^{N}\\frac{\\Theta_{j}d_{j}}{\\delta x^{p_{j}}}% \\right)^{-1}\\frac{\\delta t\\epsilon\\Theta_{i}}{\\delta x^{p_{i}}}$ as $\\Lambda_{i}$, and, we drop the subscript of $u$ and use a superscript to denote a single iteration at a time $t+\\delta t$. We enforce Dirichlet boundary condition using a projection step with a binary boundary mask $G$.\n\n $u^{m+1}=G\\left(\\sum_{i=1}^{N}\\Lambda_{i}(\\partial_{i}-d_{i}I)u^{m}+c\\right)+(I% -G)b$ (4)\n\nEq. 4 can be seen as a linear operator $u^{m+1}=\\Psi(u^{m})=Tu^{m}+k$. We can guarantee the spectral radius of the linear transformer $T$, i.e. $\\rho(T)<1$, by appropriately selecting $\\delta x,\\delta t,\\mbox{ and }\\epsilon$ [see Appendix A], leading to a fixed-point algorithm.", null, "Figure 1: Qualitative comparison of u (c.f. Eq 7) from the neural scheme (10 iterations) and a semi-implicit scheme (25 iterations) against the FEniCS solution for a test sequence of 10 time points. All methods use same initial- and boundary condition. The neural update shows consistently faster convergence than semi-implicit one.\n\n### 2.2 Neural Solver\n\nWe propose the following iterator\n\n $\\Phi_{H}(u)=\\Psi(u)+G\\left(\\sum_{i=1}^{N}\\Lambda_{i}H_{i}w\\right)$ (5)\n\nwhere $w=\\Psi(u)-u$, and $H_{i}$ is a learned linear operator which satisfies $H_{i}0=0$ for $i=1:N$. Substituting $w$ in Eq. 5 we get $\\Phi_{H}(u)=T^{\\prime}u+k^{\\prime}$, where $k^{\\prime}$ denotes the additive part which is independent of $u$, and $T^{\\prime}=T+G\\sum_{i=1}^{N}\\Lambda_{i}H_{i}(T-I)$.\n\n###### Lemma 1.\n\nFor a linear PDE problem $(\\{\\Theta_{i},\\partial_{i}\\}_{i=1:N},G,u_{0},b,\\delta x,\\delta t,\\epsilon)$ and choice of $\\{H_{i}\\}_{i=1:N}$ if $u^{*}$ is a fixed point of $\\Psi,$ it is a fixed point of $\\Phi_{H}$ in Eq. 5.\n\n###### Proof.\n\nBased on the iterative rule in Eq.5 $,$ if $u^{*}$ satisfies $\\Psi\\left(u^{*}\\right)=u^{*}$ then $w=\\Psi\\left(u^{*}\\right)-u^{*}=0$ . Therefore, $\\Phi_{H}\\left(u^{*}\\right)=\\Psi\\left(u^{*}\\right)+G\\sum_{i=1}^{N}\\Lambda_{i}H_% {i}0=u^{*}.$\n\nMoreover, the space of $\\Phi_{H}$ subsumes the standard solver $\\Psi.$ If $H_{i}=0,$ then $\\Phi_{H}=\\Psi.$ Furthermore, if $H_{i}=\\partial_{i}$ , then since $GT=T$\n\n $\\Phi_{H}(u)=\\Psi(u)+GT(\\Psi(u)-u)=T\\Psi(u)+c=\\Psi^{2}(u)$ (6)\n\nwhich is equal to two iterations of $\\Psi.$ Because computing $\\Phi$ requires two convolutions: ${H_{i}}$ and $T$ one iteration with $\\Phi_{H}$ has the same number of convolution operations as two iterations of $\\Psi.$ This shows that we can learn an ${H_{i}}$ such that our iterator $\\Phi_{H}$ is at least as the standard solver $\\Psi.$ In the following theorem, we show that there is a convex open set of ${H_{i}}$ that the learning algorithm can explore.\n\n###### Theorem 1.\n\nFor fixed $G,\\{\\Theta_{i}\\}_{i=1:N},u_{0},b,\\delta x,\\delta t$, and $\\epsilon$ the spectral norm of $\\Phi_{H}(u;G,\\{\\Theta_{i}\\}_{i=1:N},u_{0},b,\\delta x,\\delta t,\\epsilon)$ is a convex function of $\\{H_{i}\\}_{i=1:N}$ and the set of $\\{H_{i}\\}_{i=1:N}$ such that the spectral norm of $\\Phi_{H}(u)<1$ is a convex open set.\n\n###### Proof.\n\nSee Appendix B. ∎\n\nOn a stark contrast with previous work , we have several sets of parameters $\\{\\Theta_{i}\\}_{i=1:N},\\delta x,\\delta t$, and $\\epsilon$ attached to the PDEs governing equation. Although we train on a single domain, the model posses a generalization properties, which we show in the following.\n\n###### Proposition 1.\n\nFor fixed {$\\partial_{i}\\}_{i=1:N},G$, and $\\{H_{i}\\}_{i=1:N}$, and some $u^{\\prime}_{0},b^{\\prime},\\{\\Theta^{\\prime}_{i}\\}_{i=1:N},\\delta x^{\\prime},% \\delta t^{\\prime}$, and $\\epsilon^{\\prime}$, if $\\Phi_{H}(u)$ is valid iterator for the PDE problem $(\\{\\Theta^{\\prime}_{i},\\partial_{i}\\}_{i=1:N},G,u^{\\prime}_{0},b^{\\prime},% \\delta x^{\\prime},\\delta t^{\\prime},\\epsilon^{\\prime})$, then for all $u_{0}$ and $b,$ the iterator $\\Phi_{H}(u)$ is valid iterator for the PDE problem $(\\{\\Theta_{i},\\partial_{i}\\}_{i=1:N},G,u_{0},b,\\delta x,\\delta t,\\epsilon)$ if $\\left\\|\\Lambda_{i}\\right\\|<\\frac{1}{\\sum_{j=1}^{N}\\left\\|\\partial_{j}-d_{j}I% \\right\\|},\\leavevmode\\nobreak\\ \\forall i=1:N$.\n\n###### Proof.\n\nSee Appendix B. ∎\n\n## 3 Experiments\n\nWe consider a 2-D advection-diffusion equation of the following form\n\n $\\displaystyle\\frac{\\partial u}{\\partial t}=\\mathbf{v}^{\\top}\\cdot\\begin{% bmatrix}\\partial_{x}u\\\\ \\partial_{y}u\\end{bmatrix}+\\mathbf{D}^{\\top}\\cdot\\begin{bmatrix}\\partial_{xx}u% \\\\ \\partial_{yy}u\\end{bmatrix}\\mbox{; subject to }u_{t_{0}}=u_{0}(x,y)$ (7)\n\nwhere $\\mathbf{v}=[v_{x},v_{y}]^{\\top}$ and $\\mathbf{D}=[D_{xx},D_{yy}]^{\\top}$ are advection velocity and diffusivity respectively. We minimize the following loss function\n\n $\\mathcal{L}=\\frac{1}{L}\\sum_{n=1}^{N}\\sum_{l=1}^{L}\\left\\|\\Phi^{n,k}_{H}(u^{l}% _{t})-u^{l}_{t+n\\delta t}\\right\\|;k\\sim\\mathcal{U}[1,20]$\n\nwhere $n$ is the number of time-step and $k$ iteration for a single time step is denoted as $\\Phi_{H}(\\Phi_{H}\\dots(\\Phi_{H}))=\\Phi^{k}_{H}$.\n\nData Generation: We consider a rectangular domain of $\\Omega=[0,2\\pi]\\times[0,2\\pi]$. Elements of $\\mathbf{v}$ and $\\mathbf{D}$ are drawn from a uniform distribution of $\\mathcal{U}[-2.0,2.0]$ and $\\mathcal{U}[0.2,0.8]$ respectively. The computational domain is discretized into 64 x 64 regular mesh. We assume zero Dirichlet boundary condition and the initial value is generated according to as $u_{0}=\\lambda cos(kx+ly)+\\gamma sin(kx+ly)$ where $\\gamma$ and $\\lambda$ are drawn from a normal distribution of $\\mathcal{N}(0,0.02)$, and, $k$ and $l$ are values drawn in random from a uniform distribution of $\\mathcal{U}[1,9]$. We generate 200 simulations each having 50 time steps, using FEniCS for $\\delta t=0.2$. Further, we take the train, test, and validation split of the simulated time series as $80\\%:10\\%:10\\%$. A time series of a test data is shown in Fig 1.\n\nImplementation Details: Following , we use a three-layer convolutional neural network to model each of the $H_{i}$. We use zero-padding to enforce zero Dirichlet condition at the boundary and use a kernel size of 3x3. During training, we fixed the following parameters as follows $\\delta x=0.098,\\delta t=0.2,\\epsilon=0.9$. We use PyTorch framework and train our network with Adam Optimizer .\n\n### 3.1 Results and Discussion", null, "Figure 2: (a), (b), and (c) shows the mean-squared error (between FEniCS solution and Semi-implicit scheme and neural scheme) vs time plot for different parameters during test time as specified. The banded curves indicate the 25% and 75% percentile of the normalized errors among 20 test samples.\n\nGeneralization for Different Parameters: We investigate the effect of different parameter settings than those used during, to validate the generalizability of the neural scheme. To study the effect of different $\\epsilon$ we use the original test set. We generate two additional test cases; varying one parameter at a time : a) $\\delta t=0.12$, and b) $\\delta x=0.049$. The elements of $\\mathbf{v}$ and $\\mathbf{D}$ are drawn from the same distribution as before. The average error propagation over ten time step between semi-implicit finite difference method and the proposed neural implicit solver is being compared in Figure 2.\n\nWe observe that error from the neural scheme (10 iterations per time step) is less compared to the error from the semi-implicit FDM (25 iterations per time step) scheme for all three different test sets. This affirms our hypothesis that the neural solver is more accurate compared to the semi-implicit FDM and generalizable to other parameter settings at the same time.\n\nRun-time Comparison: We compare the run time for the neural solver (10 iterations per time step) and semi-implicit scheme (25 iterations per time step), for one time-steps. The experiments are conducted on an Intel Xeon W-2123 CPU @ 3.60GHz, with code running on one of the four cores. We report that the trained neural solver takes circa 0.0148s compared to 0.0141s for the semi-implicit scheme, whereas the FEniCS solution takes 3.19s for machine precision convergence.\n\n## 4 Conclusion\n\nThis abstract introduces a novel implicit neural scheme to solve linear time-dependent PDEs. We leverage an existing semi-implicit update rule to design a learnable iterator which provides theoretical guarantees. The learned iterator achieves faster convergence compared to the existing semi-implicit solver and produces a more accurate solution for a fixed computation budget. More importantly, we empirically demonstrate that training on a single parameter setting is enough to generalize over other parameter settings which confirms our theoretical results. The learned neural solver can be a faster alternative for simulation of various physical processes.\n\n#### Acknowledgments\n\nSuprosanna Shit and Ivan Ezhov are supported by the Translational Brain Imaging Training Network (TRABIT) under the European Union’s ‘Horizon 2020’ research & innovation programme (Grant agreement ID: 765148).\n\n## References\n\n• M. Alnæs et al. (2015) The FEniCS project version 1.5. Archive of Numerical Software 3 (100). Cited by: §3.\n• J. Cai et al. (2012) Image restoration: total variation, wavelet frames, and beyond. Journal of the American Mathematical Society 25 (4), pp. 1033–1089. Cited by: Appendix C.\n• I. Ezhov et al. (2019) Neural parameters estimation for brain tumor growth modeling. In Proceedings of the International Conference on Medical Image Computing and Computer Assisted Intervention, Cited by: §1.\n• R. Horn and C. Johnson (1991) Topics in matrix analysis. Cambridge University Press. Cited by: Appendix A, Appendix A.\n• J. Hsieh et al. (2019) Learning neural PDE solvers with convergence guarantees. In Proceedings of the International Conference on Learning Representations, Cited by: Appendix B, §1, §1, §2.2, §3.\n• D. P. Kingma and J. Ba (2014) Adam: a method for stochastic optimization. In Proceedings of the International Conference on Learning Representations, Cited by: §3.\n• A. Logg et al. (2012) Automated solution of differential equations by the finite element method: the FEniCS book. Vol. 84, Springer Science & Business Media. Cited by: §3.\n• Z. Long et al. (2018) PDE-net: learning PDEs from data. In Proceedings of the 35th International Conference on Machine Learning, Vol. 80, pp. 3208–3216. Cited by: §1, §3.\n• M. Magill et al. (2018) Neural networks trained to solve differential equations learn general representations. In Advances in Neural Information Processing Systems, pp. 4071–4081. Cited by: §1.\n• M. Raissi et al. (2019) Physics-informed neural networks: a deep learning framework for solving forward and inverse problems involving nonlinear partial differential equations. Journal of Computational Physics 378, pp. 686 – 707. Cited by: §1.\n• J. Tompson et al. (2017) Accelerating Eulerian fluid simulation with convolutional networks. In Proceedings of the 34th International Conference on Machine Learning, Vol. 70, pp. 3424–3433. Cited by: §1.\n\n## Appendix A Convergence Criteria for Semi-implicit Update\n\nThe spectral radius of $T$ can be expresses as following\n\n $\\displaystyle\\rho(T)$ $\\displaystyle\\leq\\left\\|T\\right\\|$ $\\displaystyle=\\left\\|G\\sum_{i=1}^{N}\\Lambda_{i}(\\partial_{i}-d_{i}I)\\right\\|$ $\\displaystyle\\leq\\left\\|G\\right\\|\\sum_{i=1}^{N}\\left\\|\\Lambda_{i}\\right\\|\\left% \\|(\\partial_{i}-d_{i}I)\\right\\|;\\leavevmode\\nobreak\\ [\\mbox{ Invoking norm % inequalities\\@@cite[cite]{[\\@@bibref{}{horn1991topics}{}{}]}}]$ $\\displaystyle=\\sum_{i=1}^{N}\\left\\|\\Lambda_{i}\\right\\|\\left\\|(\\partial_{i}-d_{% i}I)\\right\\|;\\leavevmode\\nobreak\\ [\\|G\\|=1]$\n\nGiven $\\left\\|\\Lambda_{i}\\right\\|<\\frac{1}{\\sum_{j=1}^{N}\\left\\|(\\partial_{j}-d_{j}I)% \\right\\|};\\leavevmode\\nobreak\\ \\forall i=1:N$, we have $\\rho(T)<1$.\n\n## Appendix B Proofs\n\nSee 1\n\n###### Proof.\n\nThe spectral norm $\\|\\cdot\\|$ is convex from the sub-additive property, and $T^{\\prime}$ is linear in $\\{H_{i}\\}_{i=1:N}$. To prove that it is open, observe that $\\|\\cdot\\|$ is a continuous function, so $(T+G\\sum_{i=1}^{N}\\Lambda_{i}H_{i}(T-I))$ is continuous in $\\{H_{i}\\}_{i=1:N}$. Given $\\rho(T^{\\prime})$, the set of $\\{H_{i}\\}_{i=1:N}$ is the preimage under this continuous function of $(0,1-\\zeta)$ for some $\\zeta>0$, and the inverse image of open set $(0,1-\\zeta)$ must be open. ∎\n\n###### Lemma 2.\n\nThe upper bound of the spectral norm of $\\Phi_{H}$ is independent of $\\{\\Theta_{i}\\}_{i=1:N},\\delta x,\\delta t$, and $\\epsilon$ Given $\\left\\|\\Lambda_{i}\\right\\|<\\frac{1}{\\sum_{j=1}^{N}\\left\\|\\partial_{j}-d_{j}I% \\right\\|},\\leavevmode\\nobreak\\ \\forall i=1:N$.\n\n###### Proof.\n\nConsidering the spectral norm of $T^{\\prime}$ and invoking product and triangular inequality of norm, we obtain the following tight bound\n\n $\\displaystyle\\left\\|T^{\\prime}\\right\\|$ $\\displaystyle=\\left\\|G\\sum_{i=1}^{N}\\Lambda_{i}(\\partial_{i}-d_{i}I-H_{i})+G% \\sum_{i=1}^{N}(\\Lambda_{i}H_{i})T\\right\\|$ $\\displaystyle\\leq\\|G\\|\\sum_{i=1}^{N}\\left\\|\\Lambda_{i}\\right\\|\\left\\|\\partial_% {i}-d_{i}I-H_{i}\\right\\|+\\|G\\|\\sum_{i=1}^{N}\\left\\|\\Lambda_{i}\\right\\|\\left\\|H% _{i}\\right\\|\\left\\|T\\right\\|$ $\\displaystyle<\\sum_{i=1}^{N}\\left\\|\\Lambda_{i}\\right\\|\\left(\\left\\|\\partial_{i% }-d_{i}I-H_{i}\\right\\|+\\left\\|H_{i}\\right\\|\\right)\\leavevmode\\nobreak\\ [\\|G\\|=% 1,\\|T\\|<1]$\n\nGiven $\\left\\|\\Lambda_{i}\\right\\|<\\frac{1}{\\sum_{j=1}^{N}\\left\\|\\partial_{j}-d_{j}I% \\right\\|},\\leavevmode\\nobreak\\ \\forall i=1:N$, we have\n\n $\\displaystyle\\left\\|T^{\\prime}\\right\\|<\\frac{1}{\\sum_{j=1}^{N}\\left\\|\\partial_% {j}-d_{j}I\\right\\|}\\sum_{i=1}^{N}\\left(\\left\\|\\partial_{i}-d_{i}I-H_{i}\\right% \\|+\\left\\|H_{i}\\right\\|\\right)$\n\nSee 1\n\n###### Proof.\n\nFrom Theorem 1 of and Lemma 1, our iterator is valid if and only if $\\rho\\left(T^{\\prime}\\right)<1$. From Lemma 2 the upper bound of spectral norm of iterator only depends on {$\\partial_{i}\\}_{i=1:N}$ and $\\{H_{i}\\}_{i=1:N}$ given $\\left\\|\\Lambda_{i}\\right\\|<\\frac{1}{\\sum_{j=1}^{N}\\left\\|\\partial_{j}-d_{j}I% \\right\\|},\\leavevmode\\nobreak\\ \\forall i=1:N$. Nonetheless, for any matrix spectral radius is upper bounded by its spectral norm. Thus, if the iterator is valid for some $u^{\\prime}_{0},b^{\\prime},\\{\\Theta^{\\prime}_{i}\\}_{i=1:N},\\delta x^{\\prime},% \\delta t^{\\prime}$, and $\\epsilon^{\\prime}$, then it is valid for any feasible choice of $u_{0},b,\\{\\Theta_{i}\\}_{i=1:N},\\delta x,\\delta t$, and $\\epsilon$ satisfying the constraints. ∎\n\n## Appendix C Geometric Interpretation\n\nSurprisingly, we find that the form of $\\|T^{\\prime}\\|$ has a special structure. As the denominator is constant the objective is to minimize $\\left\\|\\partial_{i}-d_{i}I-H_{i}\\right\\|+\\left\\|H_{i}\\right\\|$ w.r.t. $\\left\\|H_{i}\\right\\|;\\forall i=1:N$. Invoking triangular inequality of norm we have the lower-bound\n\n $\\left\\|\\partial_{i}-d_{i}I-H_{i}\\right\\|+\\left\\|H_{i}\\right\\|\\geq\\left\\|% \\partial_{i}-d_{i}I\\right\\|$\n\nTaking square of both side, we have\n\n $\\mbox{Find }{H_{i}}\\mbox{ such that }(\\partial_{i}-d_{i}I)^{\\top}H_{i}\\leq% \\left\\|H_{i}\\right\\|$\n\nWhen the equality holds the local optima is the surface of the hyper-sphere with center at $\\frac{1}{2}(\\partial_{i}-d_{i}I)$ with radius $\\frac{1}{2}\\left\\|\\partial_{i}-d_{i}I\\right\\|$. We interpret that each $H_{i}$ explores an optimal discretization scheme near the manifold of $\\partial_{i}$ by learning the sum of order rules as described in ." ]
[ null, "https://deeplearn.org/arxiv_files/1910.03452v1/Figs/sequence.png", null, "https://deeplearn.org/arxiv_files/1910.03452v1/Figs/delta_0_75_var.png", null ]
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https://www.cuemath.com/8th-grade-math/
[ "By Grade 8 children are well versed with a bunch of topics. Higher-order concepts like rational numbers, irrational numbers, exponents are all introduced in Grade 8. Children in Grade 8 are expected to practice these concepts and get a hold over them. Concepts learnt in Grade 8 go a long way in higher grade concepts.\n\n## Exponents\n\n• Basics of Exponents\n• Laws of Exponents\n• Negative Exponents\n• ## Scientific Notation\n\n• Expressing Very Large and Very Small Numbers in the Standard Form\n• Operations on Numbers Expressed in Standard Form\n• ## Lines, Angles, and Triangles\n\n• Angles Formed by Intersection of Parallel Lines and Transversals\n• The Angle Sum Property and Exterior Angle Property of Triangles\n\n• Translations\n• Reflections\n• Rotations\n• ## Congruence\n\n• Congruent Figures\n• ## Similarity\n\n• Dilation\n• Similar Figures\n• ## Linear Equations in One Variable\n\n• Solving Simple Linear Equations in One Variable\n• Solving Equations Reducible to the Linear Form\n• Applications of Linear Equations in One Variable\n• ## Straight Lines\n\n• Graphs of Proportional Relationships\n• Comparing Proportional Relationships\n• The Slope-Intercept Form of a Line's Equation\n• ## Simultaneous Linear Equations\n\n• Solving a Pair of Linear Equations Graphically\n• Solving a Pair of Linear Equations Algebraically\n• Applications of Simultaneous Linear Equations\n• ## Functions\n\n• Relation and Functions\n• Representation of Functions\n• Linear and Nonlinear Functions\n• Analyzing and Sketching Graphs\n• ## Volume\n\n• Volume of Cylinders\n• Volume of Cones\n• Volume of Spheres\n• ## Data Analysis and Displays\n\n• Scatter Plots\n• Lines of Fit\n• Two-Way Tables\n• Choosing a Data Display\n• ## Real Numbers\n\n• Classification of Rational Numbers\n• Conversion of Recurring Decimals to Fractions\n• Representation of Recurring Decimals on the Number Line\n• Irrational Numbers\n• ## Square Roots and Cube Roots\n\n• Finding Square Roots\n• Finding Cube Roots\n• ## The Pythagoras Theorem\n\n• The Pythagoras Theorem and Its Application\n• The Distance Formula\n\n• All Topics\n• ## 8th Grade Math Curriculum\n\nScoring high marks in Grade 8 Mathematics just got easier! Your kid will now be familiar with Translations, Reflections, Rotations & will start performing operations on exponents with his learnings from Cuemath Online Tuitions. Your child will start asking the why behind every fact, instead of being satisfied being just the what of that fact..get ready to have a curious grade 8 kid, who just wants to be tutored more!\n\nBy the end of this grade, your kids will learn to solve simple linear equations in one variable. Cuemath tutors will encourage your kids to learn the concepts using objects, pictures, and visual models. In the online Math tuitions with Cuemath, your kid will get personalized attention right from the first class in Grade 8. Our tutors are handpicked and the best in the market. Your child will be tutored online for every concept in grade 8 by an expert teacher online. As each tuition class progresses, you will see the improvement in your kids’ mathematical abilities.\n\n## Best tutors for kids in Grade 8\n\nOur Grade 8 online tuition classes with Cuemath tutors will help your kid tackle any Math problem. Cuemath’s online tuitions are designed to teach, challenge, and boost the confidence of budding Mathematicians. One of the most important things to understand about your kid in Grade 8 is that the more he practices in online tuitions, the faster he will grasp the formulas and techniques to get ahead. Get your child enrolled in Cuemath online classes today, and see a difference right after a couple of tuitions!\n\n>> Visualize Math with Interactive 8th Grade Math Worksheets\n\n## Grade 8 kids can learn Math Online\n\nCuemath has over a dozen unique activities and alluring characters- for kids in grade 8 who are tutoring online with us. During these activities, the tutors observe Grade 8 students picking up on Relation and Functions, Analyzing and Sketching Graphs, and so many more things. Cuemath tutors observe the performance of your kid and create a personalized learning plan instead of robotic completion of courses." ]
[ null ]
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https://web2.0calc.com/questions/math-homework-help-please_1
[ "+0\n\n0\n323\n3\n\n1) The three angle measures of a triangle are x/2, x/3 and x/6 degrees. What is the measure of the smallest angle?\n\n2)Each pair of opposite sides of an octagon are parallel. We know three interior angles of the octagon are 115,125  and,135. What is the degree measure of the largest interior angle of the octagon?\n\n3)What is the ratio of the number of degrees in the interior angle of a regular pentagon to the number of degrees in the interior angle of a regular hexagon?\n\nThanks!\n\nFeb 5, 2019\n\n#1\n+1\n\n1) If x/2, x/3, and x/6 are angle measures of a triangle, they must add up to 180.\n\nx/2 + x/3 + x/6 = 180\n\nFind the common denominator:\n\n3x/6+2x/6 + x/6 = 180\n\nMultiply both sides by 6:\n\n3x+2x+x=1080\n\n6x=1080,\n\nx = 180\n\nx/6, the smallest angle, is 30.\n\n2) Correct me if I'm wrong, but isn't it just 135?\n\n3) 108/120 = 9/10\n\nFeb 5, 2019\n#2\n0\n\nFor #2, 135 isn't correct but thank you for helping me solve 1 and 3. I did not look at #1 in the way that you showed.\n\nFeb 5, 2019\n#3\n+1\n\nI believe this might be the second one\n\nThe sum of the interior angles of an octagon = 135 * 8  = 1080°\n\nSo....the remaining two angles must sum to   1080 - 2(115 + 125 + 135)  = 330°\n\nSo....since the opposite sides are parallel.....the largest remaining angles each measure 330/2 = 165°", null, "", null, "", null, "Feb 6, 2019" ]
[ null, "https://web2.0calc.com/img/emoticons/smiley-cool.gif", null, "https://web2.0calc.com/img/emoticons/smiley-cool.gif", null, "https://web2.0calc.com/img/emoticons/smiley-cool.gif", null ]
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https://math.stackexchange.com/questions/54506/is-this-batman-equation-for-real/54521
[ "# Is this Batman equation for real? [closed]\n\nHardOCP has an image with an equation which apparently draws the Batman logo. Is this for real?", null, "• Why don't you just try it? Jul 29, 2011 at 21:19\n• @Jim: If you mouse over the downvote button, you see: \"This question does not show any research effort; it is unclear or not useful.\" I downvoted because the OP was too lazy to type in the equation himself to any plotting program or calculator, which would have immediately shown that the equation is \"for real\". If the OP were asking for an explanation of how such an equation might be derived, as ShreevatsaR has done, that would be an appropriate question. Jul 30, 2011 at 17:08\n• @Zev Chonoles: i do not know of any web-site that can plot that equation. i wouldn't know where to begin. Also i don't understand how any solver could plot such a diagram. (Hence the question). Jul 30, 2011 at 19:35\n• I don't understand why this question has so many upvotes. Aug 3, 2011 at 18:36\n• Aug 31, 2011 at 14:40\n\nAs Willie Wong observed, including an expression of the form $\\displaystyle \\frac{|\\alpha|}{\\alpha}$ is a way of ensuring that $\\alpha > 0$. (As $\\sqrt{|\\alpha|/\\alpha}$ is $1$ if $\\alpha > 0$ and non-real if $\\alpha < 0$.)\n\nThe ellipse $\\displaystyle \\left( \\frac{x}{7} \\right)^{2} + \\left( \\frac{y}{3} \\right)^{2} - 1 = 0$ looks like this:", null, "So the curve $\\left( \\frac{x}{7} \\right)^{2}\\sqrt{\\frac{\\left| \\left| x \\right|-3 \\right|}{\\left| x \\right|-3}} + \\left( \\frac{y}{3} \\right)^{2}\\sqrt{\\frac{\\left| y+3\\frac{\\sqrt{33}}{7} \\right|}{y+3\\frac{\\sqrt{33}}{7}}} - 1 = 0$ is the above ellipse, in the region where $|x|>3$ and $y > -3\\sqrt{33}/7$:", null, "That's the first factor.\n\nThe second factor is quite ingeniously done. The curve $\\left| \\frac{x}{2} \\right|\\; -\\; \\frac{\\left( 3\\sqrt{33}-7 \\right)}{112}x^{2}\\; -\\; 3\\; +\\; \\sqrt{1-\\left( \\left| \\left| x \\right|-2 \\right|-1 \\right)^{2}}-y=0$ looks like:", null, "This is got by adding $y = \\left| \\frac{x}{2} \\right| - \\frac{\\left( 3\\sqrt{33}-7 \\right)}{112}x^{2} - 3$, a parabola on the positive-x side, reflected:", null, "and $y = \\sqrt{1-\\left( \\left| \\left| x \\right|-2 \\right|-1 \\right)^{2}}$, the upper halves of the four circles $\\left( \\left| \\left| x \\right|-2 \\right|-1 \\right)^2 + y^2 = 1$:", null, "The third factor $9\\sqrt{\\frac{\\left( \\left| \\left( 1-\\left| x \\right| \\right)\\left( \\left| x \\right|-.75 \\right) \\right| \\right)}{\\left( 1-\\left| x \\right| \\right)\\left( \\left| x \\right|-.75 \\right)}}\\; -\\; 8\\left| x \\right|\\; -\\; y\\; =\\; 0$ is just the pair of lines y = 9 - 8|x|:", null, "truncated to the region $0.75 < |x| < 1$.\n\nSimilarly, the fourth factor $3\\left| x \\right|\\; +\\; .75\\sqrt{\\left( \\frac{\\left| \\left( .75-\\left| x \\right| \\right)\\left( \\left| x \\right|-.5 \\right) \\right|}{\\left( .75-\\left| x \\right| \\right)\\left( \\left| x \\right|-.5 \\right)} \\right)}\\; -\\; y\\; =\\; 0$ is the pair of lines $y = 3|x| + 0.75$:", null, "truncated to the region $0.5 < |x| < 0.75$.\n\nThe fifth factor $2.25\\sqrt{\\frac{\\left| \\left( .5-x \\right)\\left( x+.5 \\right) \\right|}{\\left( .5-x \\right)\\left( x+.5 \\right)}}\\; -\\; y\\; =\\; 0$ is the line $y = 2.25$ truncated to $-0.5 < x < 0.5$.\n\nFinally, $\\frac{6\\sqrt{10}}{7}\\; +\\; \\left( 1.5\\; -\\; .5\\left| x \\right| \\right)\\; -\\; \\frac{\\left( 6\\sqrt{10} \\right)}{14}\\sqrt{4-\\left( \\left| x \\right|-1 \\right)^{2}}\\; -\\; y\\; =\\; 0$ looks like:", null, "so the sixth factor $\\frac{6\\sqrt{10}}{7}\\; +\\; \\left( 1.5\\; -\\; .5\\left| x \\right| \\right)\\sqrt{\\frac{\\left| \\left| x \\right|-1 \\right|}{\\left| x \\right|-1}}\\; -\\; \\frac{\\left( 6\\sqrt{10} \\right)}{14}\\sqrt{4-\\left( \\left| x \\right|-1 \\right)^{2}}\\; -\\; y\\; =\\; 0$ looks like", null, "As a product of factors is $0$ iff any one of them is $0$, multiplying these six factors puts the curves together, giving: (the software, Grapher.app, chokes a bit on the third factor, and entirely on the fourth)", null, "• I tip my hat to you for this comprehensive dissection. Jul 30, 2011 at 14:06\n• If there were only no rep-cap, Shree would be swimming in rep now... ;P Jul 30, 2011 at 19:56\n• \"I don’t know how ShreevatsaR did it but he sure brought a large luggage when intelligence showered the Earth.\" yangkidudel.wordpress.com/2011/08/02/love-and-mathematics Aug 3, 2011 at 0:15\n• @Jonas Meyer: LOL, that's embarrassing! :P But then again, if Batman is what it takes for someone to appreciate mathematics a little, well good for Batman. :-) Aug 3, 2011 at 10:46\n• @Jack: Grapher, which comes by default on Mac OS X. I mentioned it in the answer actually, just before the last figure. Aug 3, 2011 at 17:34\n\nYou may be able to see more easily the correspondences between the equations and the graph through the following picture which is from the link I got after a curious search on Google(link broken now):", null, "• Geometer's Sketchpad? Aug 3, 2011 at 18:28\n• @Isaac: Probably right. Sep 6, 2011 at 12:27\n• to see the graph just google the equation: 2*sqrt(-abs(abs(x)-1)*abs(3-abs(x))/((abs(x)-1)*(3-abs(x))))(1+abs(abs(x)-3)/(abs(x)-3))sqrt(1-(x/7)^2)+(5+0.97(abs(x-.5)+abs(x+.5))-3(abs(x-.75)+abs(x+.75)))(1+abs(1-abs(x))/(1-abs(x))),-3sqrt(1-(x/7)^2)sqrt(abs(abs(x)-4)/(abs(x)-4)),abs(x/2)-0.0913722(x^2)-3+sqrt(1-(abs(abs(x)-2)-1)^2),(2.71052+(1.5-.5abs(x))-1.35526sqrt(4-(abs(x)-1)^2))sqrt(abs(abs(x)-1)/(abs(x)-1))+0.9 Aug 20, 2012 at 18:47\n• @HelderVelez 2014: \"'abs' (and any subsequent words) was ignored because we limit queries to 32 words.\" - Google\n– Baby\nNov 28, 2014 at 8:11\n• x1(y), x2(y) should be 7 * sqrt( 1 - y^2 / 9) and -7 * sqrt( 1 - y^2 / 9), respectively Jun 25, 2017 at 4:39\n\nHere's what I got from the equation using Maple...", null, "• What if Commissioner Gordon uses Mathematica ???? Jul 30, 2011 at 5:25\n• @The Chaz: Then Commissioner Gordon should support the Mathematica SE Site Proposal on Area 51. Jul 30, 2011 at 19:22\n• touché .... Jul 30, 2011 at 19:24\n• Here's Mathematica code. See Heike's post. I tried it on M8 and it works fine. groups.google.com/group/comp.soft-sys.math.mathematica/…\n– Sol\nAug 5, 2011 at 22:48\n• @Sol: You can do better; see my answer. Aug 9, 2011 at 7:56\n\nLooking at the equation, it looks like it contains terms of the form $$\\sqrt{\\frac{| |x| - 1 |}{|x| - 1}}$$ which evaluates to $$\\begin{cases} 1 & |x| > 1\\\\ i & |x| < 1\\end{cases}$$\n\nSince any non-zero real number $y$ cannot be equal to a purely imaginary non-zero number, the presence of that term is a way of writing a piece-wise defined function as a single expression. My guess is that if you try to plot this in $\\mathbb{C}^2$ instead of $\\mathbb{R}^2$ you will get all kinds of awful.\n\n• Yeah, the equation looks too contrived to me. :) A parametric form (it's just quadratic and linear arcs sewn together, it looks) would still be messy, but not as messy. (Probably a good job for splines...) Jul 30, 2011 at 2:43\n• +1 i was wondering how they split it up into sections. Jul 30, 2011 at 19:38\n• \" My guess is that if you try to plot this in C2 instead of R2 you will get all kinds of awful.\" What did you expect? The analytic continuation of the Batman symbol??\n– jwg\nAug 1, 2013 at 8:43\n\nSince people (not from this site, but still...) keep bugging me, and I am unable to edit my previous answer, here's Mathematica code for plotting this monster:\n\nPlot[{With[{w = 3 Sqrt[1 - (x/7)^2],\nl = 6/7 Sqrt + (3 + x)/2 - 3/7 Sqrt Sqrt[4 - (x + 1)^2],\nh = (3 (Abs[x - 1/2] + Abs[x + 1/2] + 6) -\n11 (Abs[x - 3/4] + Abs[x + 3/4]))/2,\nr = 6/7 Sqrt + (3 - x)/2 - 3/7 Sqrt Sqrt[4 - (x - 1)^2]},\nw + (l - w) UnitStep[x + 3] + (h - l) UnitStep[x + 1] +\n(r - h) UnitStep[x - 1] + (w - r) UnitStep[x - 3]],\n1/2 (3 Sqrt[1 - (x/7)^2] + Sqrt[1 - (Abs[Abs[x] - 2] - 1)^2] + Abs[x/2] -\n((3 Sqrt - 7)/112) x^2 - 3) (Sign[x + 4] - Sign[x - 4]) - 3*Sqrt[1 - (x/7)^2]},\n{x, -7, 7}, AspectRatio -> Automatic, Axes -> None, Frame -> True,\nPlotStyle -> Black]", null, "This should work even for versions that do not have the Piecewise[] construct. Enjoy. :P\n\nIn fact, the five linear pieces that consist the \"head\" (corresponding to the third, fourth, and fifth pieces in Shreevatsa's answer) can be expressed in a less complicated manner, like so:\n\n$$y=\\frac{\\sqrt{\\mathrm{sign}(1-|x|)}}{2}\\left(3\\left(\\left|x-\\frac12\\right|+\\left|x+\\frac12\\right|+6\\right)-11\\left(\\left|x-\\frac34\\right|+\\left|x+\\frac34\\right|\\right)\\right)$$\n\nThis can be derived by noting that the functions\n\n$$\\begin{cases}f(x)&\\text{if }x<c\\\\g(x)&\\text{if }c<x\\end{cases}$$\n\nand $f(x)+(g(x)-f(x))U(x-c)$ (where $U(x)$ is the unit step function) are equivalent, and using the \"relation\"\n\n$$U(x)=\\frac{x+|x|}{2x}$$\n\nNote that the elliptic sections (both ends of the \"wings\", corresponding to the first piece in Shreevatsa's answer) were cut along the lines $y=-\\frac37\\left((2\\sqrt{10}+\\sqrt{33})|x|-8\\sqrt{10}-3\\sqrt{33}\\right)$, so the elliptic potion can alternatively be expressed as\n\n$$\\left(\\left(\\frac{x}{7}\\right)^2+\\left(\\frac{y}{3}\\right)^2-1\\right)\\sqrt{\\mathrm{sign}\\left(y+\\frac37\\left((2\\sqrt{10}+\\sqrt{33})|x|-8\\sqrt{10}-3\\sqrt{33}\\right)\\right)}=0$$\n\nTheoretically, since all you have are arcs of linear and quadratic curves, the chimera can be expressed parametrically using rational B-splines, but I'll leave that for someone else to explore...\n\n• Okay, make that an hour. Sheesh. *facepalm* Jul 30, 2011 at 19:56\n• Great, this would be much clearer. BTW, I think it's safer to leave it at \"can be expressed using B-splines\" than to actually try it out: who knows how many hours that will waste, right? :-) Jul 31, 2011 at 3:20\n\nThe following is what I got from the equations using MATLAB:", null, "Here is the M-File (thanks to this link):\n\nclf; clc; clear all;\nsyms x y\n\neq1 = ((x/7)^2*sqrt(abs(abs(x)-3)/(abs(x)-3))+(y/3)^2*sqrt(abs(y+3/7*sqrt(33))/(y+3/7*sqrt(33)))-1);\neq2 = (abs(x/2)-((3*sqrt(33)-7)/112)*x^2-3+sqrt(1-(abs(abs(x)-2)-1)^2)-y);\neq3 = (9*sqrt(abs((abs(x)-1)*(abs(x)-.75))/((1-abs(x))*(abs(x)-.75)))-8*abs(x)-y);\neq4 = (3*abs(x)+.75*sqrt(abs((abs(x)-.75)*(abs(x)-.5))/((.75-abs(x))*(abs(x)-.5)))-y);\neq5 = (2.25*sqrt(abs((x-.5)*(x+.5))/((.5-x)*(.5+x)))-y);\neq6 = (6*sqrt(10)/7+(1.5-.5*abs(x))*sqrt(abs(abs(x)-1)/(abs(x)-1))-(6*sqrt(10)/14)*sqrt(4-(abs(x)-1)^2)-y);\n\naxes('Xlim', [-7.25 7.25], 'Ylim', [-5 5]);\nhold on\n\nezplot(eq1,[-8 8 -3*sqrt(33)/7 6-4*sqrt(33)/7]);\nezplot(eq2,[-4 4]);\nezplot(eq3,[-1 -0.75 -5 5]);\nezplot(eq3,[0.75 1 -5 5]);\nezplot(eq4,[-0.75 0.75 2.25 5]);\nezplot(eq5,[-0.5 0.5 -5 5]);\nezplot(eq6,[-3 -1 -5 5]);\nezplot(eq6,[1 3 -5 5]);\ncolormap([0 0 1])\n\ntitle('Batman');\nxlabel('');\nylabel('');\nhold off\n\n\nThe 'Batman equation' above relies on an artifact of the plotting software used which blithely ignores the fact that the value $\\sqrt{\\frac{|x|}{x}}$ is undefined when $x=0$. Indeed, since we’re dealing with real numbers, this value is really only defined when $x>0$. It seems a little ‘sneaky’ to rely on the solver to ignore complex values and also to conveniently ignore undefined values.\n\nA nicer solution would be one that is unequivocally defined everywhere (in the real, as opposed to complex, world). Furthermore, a nice solution would be ‘robust’ in that small variations (such as those arising from, say, roundoff) would perturb the solution slightly (as opposed to eliminating large chunks).\n\nTry the following in Maxima (actually wxmaxima) which is free. The resulting plot is not quite as nice as the plot above (the lines around the head don’t have that nice ‘straight line’ look), but seems more ‘legitimate’ to me (in that any reasonable solver should plot a similar shape). Please excuse the code mess.\n\n/* [wxMaxima batch file version 1] [ DO NOT EDIT BY HAND! ]*/\n/* [ Created with wxMaxima version 0.8.5 ] */\n\n/* [wxMaxima: input start ] */\n/* [wxMaxima: input end ] */\n\n/* [wxMaxima: input start ] */\nf(a,b,x,y):=a*x^2+b*y^2;\n/* [wxMaxima: input end ] */\n\n/* [wxMaxima: input start ] */\nc1:sqrt(26);\n/* [wxMaxima: input end ] */\n\n/* [wxMaxima: input start ] */\ndraw2d(implicit(\nf(1/36,1/9,x,y)\n+max(0,2-f(1.5,1,x+3,y+2.7))\n+max(0,2-f(1.5,1,x-3,y+2.7))\n+max(0,2-f(1.9,1/1.7,(5*(x+1)+(y+3.5))/c1,(-(x+1)+5*(y+3.5))/c1))\n+max(0,2-f(1.9,1/1.7,(5*(x-1)-(y+3.5))/c1,((x-1)+5*(y+3.5))/c1))\n+max(0,2-((1.1*(x-2))^4-(y-2.1)))\n+max(0,2-((1.1*(x+2))^4-(y-2.1)))\n+max(0,2-((1.5*x)^8-(y-3.5)))\n-1,\nx,-6,6,y,-4,4));\n/* [wxMaxima: input end ] */\n\n/* Maxima can't load/batch files which end with a comment! */\n\"Created with wxMaxima\"$ The resulting plot is:", null, "(Note that this is, more or less, a copy of the entry I made on http://blog.makezine.com.) • I really think that an indeterminate value multiplied by zero equals zero, so it seems to be legit. Is there any reason 0 * 0/0 should not be defined to be zero? Jun 15, 2015 at 19:45 • @dbanet: What are you referring to? The issue above is that the original equations rely on the plotting software ignoring undefined values, which is peculiar, to say the least. The expression$\\sqrt{\\frac{|x|}{x}}$(with$x$being replaced by some expression) is what I referred to and it appears without being multiplied by$x$. Jun 15, 2015 at 19:54 • @copper-hat: The function$f(x)=\\sqrt{\\frac{|x|}{x}}$appears only in boolean expressions$F:x\\to \\{\\text{True},\\text{False}\\}$of form$f(x)g(x)=0$, so if$g(x)$is defined as$g:x\\in\\mathbb{C}\\to{0}$, I would rather evaluate$F(0)$to$\\text{True}$than to$\\text{False}$, as$\\forall{x}:f(x)g(x)=0\\Longleftrightarrow \\Big(f(x)=0\\lor g(x)=0\\Big)$. Jun 15, 2015 at 21:57 • @dbanet: I'm really not sure what you are getting at. Look at the expressions in the question. They rely on the expression$\\sqrt{\\frac{|x|}{x}}$returning zero for$x \\le 0$, which is strange (look at Willie's answer math.stackexchange.com/a/54521/27978). My answer plots level sets, which avoids this whole issue. Jun 15, 2015 at 22:10 • @dbanet: I don't really get your point. In the formula in the question there are lots of expressions of the above form that are multiplied by quantities that do not evaluate to zero when$x<0\\$. If they did, there would be no need to have the strange expression in the first place. Jun 15, 2015 at 22:21\n\nHere's the equations typed out if you want save time with writing it yourself.\n\n(x/7)^2*SQRT(ABS(ABS(x)-3)/(ABS(x)-3))+(y/3)^2\\*SQRT(ABS(y+3*SQRT(33)/7)/(y+3*SQRT(33)/7))-1=0\nABS(x/2)-((3*SQRT(33)-7)/112)*x^2-3+SQRT(1-(ABS(ABS(x)-2)-1)^2)-y=0\n9*SQRT(ABS((ABS(x)-1)*(ABS(x)-0.75))/((1-ABS(x))*(ABS(x)-0.75)))-8*ABS(x)-y=0\n3*ABS(x)+0.75*SQRT(ABS((ABS(x)-0.75)*(ABS(x)-0.5))/((0.75-ABS(x))*(ABS(x)-0.5)))-y=0\n2.25*SQRT(ABS((x-0.5)*(x+0.5))/((0.5-x)*(0.5+x)))-y=0\n(6*SQRT(10))/7+(1.5-0.5*ABS(x))*SQRT(ABS(ABS(x)-1)/(ABS(x)-1))-((6*SQRT(10))/14)*SQRT(4-(ABS(x)-1)^2)-y=0\n\n• The multiple \"=0\" lines are different than the mulitiplications in the original, there are some backslashes in there that throw things off, and the formatting is hard to read. The pastebin is better. Jul 30, 2011 at 16:42\n• @copper, it is just because of the algorithm that draw uses for drawing implicit functions. You need to setup the variables ip_grid and ip_grid_in, that are the sampling values in your region. For example draw2d(ip_grid=[60,60], ip_grid_in=[20,20], implicit(y^2=x^3-2*x+1, x, -4,4, y, -4,4) ); Oct 16, 2014 at 19:21" ]
[ null, "https://i.stack.imgur.com/VYKfg.jpg", null, "https://i.stack.imgur.com/PXv4W.png", null, "https://i.stack.imgur.com/oeCdG.png", null, "https://i.stack.imgur.com/vAoFe.png", null, "https://i.stack.imgur.com/Vfyre.png", null, "https://i.stack.imgur.com/69Pdf.png", null, "https://i.stack.imgur.com/3CJGO.png", null, "https://i.stack.imgur.com/Sh0Bp.png", null, "https://i.stack.imgur.com/XKs3Z.png", null, "https://i.stack.imgur.com/OO3np.png", null, "https://i.stack.imgur.com/YmA4v.png", null, "https://i.stack.imgur.com/EHNR8.png", null, "https://i.stack.imgur.com/6pXe9.jpg", null, "https://i.stack.imgur.com/zYXKB.png", null, "https://i.stack.imgur.com/vHI8K.jpg", null, "https://i.stack.imgur.com/dMzTL.png", null ]
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https://www.esaral.com/q/two-angles-of-a-triangle-are-equal-and-the-third-angle-is-greater-than-each-of-those-angles-by-30-84307
[ "# Two angles of a triangle are equal and the third angle is greater than each of those angles by 30°.\n\nQuestion:\n\nTwo angles of a triangle are equal and the third angle is greater than each of those angles by 30°. Determine all the angles of the triangle.\n\nSolution:\n\nGiven that,\n\nTwo angles of a triangle are equal and the third angle is greater than each of those angles by 30°.\n\nLet x, x, x + 30° be the angles of a triangle\n\nWe know that,\n\nSum of all angles in a triangle is 180°\n\nx + x + x + 30° = 180°\n\n3x + 30° = 180°\n\n3x = 180° − 30°\n\n3x = 150°\n\nx = 50°\n\nTherefore, the three angles are 50°, 50°, 80°." ]
[ null ]
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https://www.programcreek.com/2014/03/create-java-string-by-double-quotes-vs-by-constructor/
[ "# Create Java String Using ” ” or Constructor?\n\nIn Java, a string can be created in two ways:\n\n```String x = \"abc\"; String y = new String(\"abc\");```\n\nWhat is the difference between using the double quotes and using the constructor?\n\n1. Double Quotes vs. Constructor\n\nThis question can be answered by using two simple examples.\n\nExample 1:\n\n```String a = \"abcd\"; String b = \"abcd\"; System.out.println(a == b); // True System.out.println(a.equals(b)); // True```\n\n`a==b` is true because `a` and `b` are referring to the same string literal in the method area. The memory references are the same.\n\nWhen the same string literal is created more than once, only one copy of each distinct string value is stored. This is called \"string interning\". All compile-time constant strings in Java are automatically interned.\n\nExample 2:\n\n```String c = new String(\"abcd\"); String d = new String(\"abcd\"); System.out.println(c == d); // False System.out.println(c.equals(d)); // True```\n\n`c==d` is false because `c` and `d` refer to two different objects in the heap. Different objects always have different memory references.\n\nThis diagram illustrate the two situations above:", null, "2. Run-Time String Interning\n\nThanks to LukasEder (his comment below):\n\nString interning can still be done at run-time, even if two strings are constructed with constructors:\n\n```String c = new String(\"abcd\").intern(); String d = new String(\"abcd\").intern(); System.out.println(c == d); // Now true System.out.println(c.equals(d)); // True```\n\n3. When to Use Which\n\nBecause the literal \"abcd\" is already of type String, using constructor will create an extra unnecessary object. Therefore, double quotes should be used if you just need to create a String.\n\nIf you do need to create a new object in the heap, constructor should be used. Here is a use case.\n\nCategory >> Basics >> Diagram >> Java\nIf you want someone to read your code, please put the code inside <pre><code> and </code></pre> tags. For example:\n```<pre><code>\nString foo = \"bar\";\n</code></pre>\n```\n• sss\n\nin jdk 1.6 is also true\n\n• affan\n\nSir I am newly joined Java programming certificate course my preceptor give me an project & I am stack on it for instance sir kindly advice me that how could I will manage program where I using string,double,integer variable in simultaneously please give some example.\n\n• Walt Corey\n\nThis is how one writes code nobody can understand. Good way to be top of the list come next layoff. I knew a guy did the same crap wit C++, so ‘highly optimized” nobody could read it, and it also didn’t work but broke in mysterious ways. Rather than seeing how obtuse you can code, see how readable you can. Readability is FAR more cost effective than so-called highly optimized code.\n\n• Ritter Liu\n\nThank you for your share 😛\n\n• Charly\n\nYou are worng, In JDK1.6 and JDK1.7 return true. Test rourself\n\n• 智 陶\n\nString c = new String(“abcd”).intern();\nString d = new String(“abcd”).intern();\nSystem.out.println(c == d); // True in JDK1.7 while false in JDK 1.6.\ncorrect me if i am wrong\n\n• Feng Sun\n\ngood to know, thanks\n\n• LukasEder\n\nWhen the same string literal is created more than once, JVM store only one copy of each distinct string value. This is called “string interning“.\n\nThat’s not entirely correct. In your example, the compiler will already make sure that the same string constant is effectively referenced. String interning can still be done at runtime, even if two strings are constructed with constructors:\n\nString c = new String(“abcd”).intern();\nString d = new String(“abcd”).intern();\nSystem.out.println(c == d); // Now true\nSystem.out.println(c.equals(d)); // True" ]
[ null, "https://www.programcreek.com/wp-content/uploads/2014/03/constructor-vs-double-quotes-Java-String-New-Page-650x324.png", null ]
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https://present5.com/ministry-of-education-and-science-of-ukraine-national/
[ "", null, "Скачать презентацию Ministry of Education and Science of Ukraine National\n\ned_-_3.ppt\n\n• Размер: 14.4 Mегабайта\n• Количество слайдов: 81\n\n## Описание презентации Ministry of Education and Science of Ukraine National по слайдам", null, "Ministry of Education and Science of Ukraine National Aerospace University n. a. N. Y. . Zhukovsky « « Kh. AI » » Department of Graphical and Computer Modeling DESCRIPTIVE GEOMETRY", null, "Surfaces.", null, "Surfaces. Classification. Determinant. Outline.", null, "Surfaces. Classification. Determinant. Outline. Fundamentals : Surface of geometric solid is multiple of boundary points of a given body. B oundary point is a center of a sphere which always contains inner and outer points for this body while its radius tend to zero. To specify a surface in a drawing means to indicate the conditions enabling us to construct each point of this surface. Representation of a surface : 1. Analytical method ; 2. Graphical method. a ) framing; b ) kinematic. Kinematic surface represents a locus of lines moving in space according to a certain law. These lines, which may be straight or curved, are called generating lines ( generatrices ). In general, there is a variety of laws for generating a certain surface. It is desirable to select those laws and shapes of the generating lines which are most simple and convenient both for representing the required surface and solving the problems associated with it. Determinant of kinematic surface is minimum number of parameters and constraints which univocally represent given surface. Determinant structure : Determinant = {(geometrical parameters )( algorithm of constructions ) }. Surface on orthographical view can be represented by projections of its determinant.", null, "Surfaces. Classification. Surfaces Curved Polyhedral (multifaceted) Kinematical Framing. R uled (linear) N onlinear Prismatic. Pyramidal. Continuous / Discrete D evelopable C onical C ylindrical Torse G eneral view C hannel surface N ondevelopable G eneral view C atalan C yclic R evolution General view generatrix Parallel shift with Plane Director (parallelism ) Plane parallel to axes. To systematize properties of surfaces we can classify them by the type of generatrices and character of their motion.", null, "Surfaces. Classification. Determinant. Outline of a Surface. For a surface to be specified, it is sufficient to have the projections of its directrix (guide line) and adequate information on the method for constructing the generatrix passing through any point of the directrix. But if it is desirable to make the representation more obvious and expressive, then it is advisable to draw also the outline of the surface, several positions of the generatrix, most important lines and points on the surface, etc. Representation of a surface on a drawing by projections of its determinant can’t provide proper obviousness. It’s reasonable to supplement drawing of a surface by its outline. Example: a Cone, represented by its determinant : КВ = {(ℓ‚ t ) (ℓ∩ t , ℓ = R ( t )}, Difficult to imagine such a figure without additional explanations. t 2 l 1 t", null, "Surfaces. Classification. Determinant. Outline of a Surface. П 2 S Contour line Outline. Projectors O 1 O 2 O 3 П 1 On parallel projection of a surface Ф on the plane П some projectors touch with the surface Ф and form tangent cylindrical projecting surface Q. Tangent curve of Q and Ф surfaces is called Contour Line (Contour Curve), its projection on the plane П – Outline of the surface Ф. This curve can be plane or space. Outline of a surface is projection of its contour line. Thus Outline is visibility border for points which belong to a surface. For instance, the sphere on the planes П 1 and П 2 has contour curves – circles q 2 and m 1. They project on the other planes as curves of visibility q 1 and m 2 – outline curves. m 2 m 1 q 1 q", null, "Surfaces. Classification. Determinant. Outline of a surface. Frontal Visible Contour. Frontal Outline Horizontal Visible Contour Horizontal Outline", null, "Surfaces of Revolution. Forming. Definitions. Parallel Meridian Generatrix Neck Axis of revolution l 1 i m 3 m 2 m 1 o 2 o 1 o 3 i Meridian plane. A surface of revolution is generated by the revolution of a curved-line or straight-line generatrix about a fixed straight line called the axis of the revolution. Determinant of the surface of revolution — generatrix and axis of revolution : SR = {( l , I )( l = R ( I ))} Each point of the generatrix describes a circle. Hence, a plane perpendicular to the axis of revolution cuts this surface in circles. Such circles are called parallels. The biggest parallel is called the equator , the smallest parallel is called the neck of a surface. Neck and equator project to outlines. A plane passing through the axis of a surface of revolution is termed a meridian plane. The line where a meridian plane intersects a surface of revolution is called the meridian of the surface. Equator Cylinder of revolution. Generatrix – l — straight line axis of revolution ί ; n — curve , equally spaced from axis of revolution. If we imagine the totality of rectilinear generatrices, revolving around the axes, and the totality of generating circles, moving along that axes, then each line of one totality (of one «family» of lines) will intersect all lines of the other totality (of the other «family» of lines). As a result, the mesh of the given surface is obtained. Other surfaces may be thought of in the same way.", null, "Surfaces of Revolution. General.", null, "Surfaces of Revolution. Forming. Surface of revolution. Generatrix – Straight Line. Relative position of Generatrix and Axes of revolution. CYLINDER CONE SHEET HYPERBOLOID OF REVOLUTION PARALLEL INTERSECTING SKEW HYPERBOL", null, "Surfaces of Revolution. Forming. Surface of revolution. Generatrix – Circle. R>r R<r R=0 i n Rr i TORUS CLOSED TORUS SPHER", null, "Surfaces of Revolution. Examples. C O B A 12′ I I=I ’10» 1′ 2′ 3′ 4′ 5′ 6’9 ‘ ‘ 11» 12» 1» 2» 3» I=I ‘ m 1 m 2 I ‘I 1’12’ 2’11’ 10′ 9′ 8′ 3′ 7′ 6′ 5′ 4’5» 4» 3» 2» 1»6» 10»9» 8» 7» 12»11»For the hyperboloid of revolution the meridian is a hyperbola. If the hyperbola is revolved about its imaginary axis, then we have a hyperboloid of revolution of one sheet. A hyperboloid of revolution of one sheet can also be generated by the revolution of a straight line if the generatrix and the axis of revolution are skew lines.", null, "Surfaces of Revolution. Examples. C O B I I ‘ A B 1 B 2 A 1 B 1 A 2 B 2 M 2 t 1 t 2 C 2 I=I ‘A 1 C 1 M 1 Neck", null, "Surfaces of Revolution. Sections by Planes. O 14 23 2 2 21 2 (1) 1 =3 1 (2) 1 =4 1 1 1 8 1 2 1 3 1 4 1 5 16 17 11 2 2 2 =8 2 5 2 4 2 =6 2 3 2 =7 21′ 2’8′ 3’7′ 4’6′ 5′ O 1 Minor axes Major axes Minor axes", null, "Surfaces of Revolution. Sections by Planes. S ELLIPSE – TWO OUTLINES FROM ONE SIDE OF AN APEX ARE INTERSECTED PARABOLA – PARALLEL TO THE OUTLINESTRAIGHT LINES – THROUGH APEX ALONG GENERATRIX HYPERBOLA – TWO OUTLINES FROM DIFFERENT SIDES OF AN APEX ARE INTERSECTED CIRCLE – PERPENDICULAR TO AN AXESCone of revolution. Conical Sections.", null, "Surfaces of Revolution. Sections by Planes. Point Ellipse Circle Double Parabola Two straight Lines Hyperbola", null, "Surfaces of Revolution. Sections by Planes. A C B D D ‘ B ‘ C ‘ O ‘ A ‘ B 2 =D 2 =O 2 C 2 C ‘A 2 C ‘D ‘ O ‘O ELLIPS", null, "Surfaces of Revolution. Sections by Planes. O ‘M 2 =N 2=O 2 M ‘K 2 J 2 =L 2 L ‘ J ‘ N ‘K ‘ L’R jl T 2 = P 2 T ‘σ 2 σ ‘ 2 T ‘ P ‘HYPERBOL", null, "Surfaces of Revolution. Sections by Planes. O ‘ M 2 =N 2=O 2 M ‘M ‘ K 2 J 2 =L 2 L ‘J ‘ N ‘K ‘ L’ R jl T 2 = P 2 T ‘ σ 2 σ ‘ 2 T ‘P ‘ PARABOL", null, "Sections of Solids. 1. Designate all character points of intersection between cutting plane and surfaces of figure. 2. Locate axes of section 1 – 7 parallel to a cutting line. Mark distances between points in question. 3. Construct true sizes of full sections from corresponding marked points ( for example, А-А section – ellipse from cylinder cutting between points 1 – 7 , rectangle from pyramid cutting between points 3 – 6 , etc. ). 4. Designate obtained section according to a standard (ISO, GOST, DSTU, etc. ). BA A B — B 2 2′ 71 2 3 4 5 6 7 7 6 5 4 1 22 1 4 3 6′ 3’1′ 3′ 6 ‘ 1 ‘ B 34 4’5 5′ A — A 3 3’6 6’ 1 Combined figure sectioned by a plane", null, "Intersection between a Solid and a Plane. Three ways how to construct conical section. 1. Using Major and Minor Axis of Ellipse 2. Using Parallels 3. Using Generatrices AA S 1 1 1 3 11 2 (2) 2 3 2 (4) 2 2 1 4 1 5 15 2 S 2 BB S 11 2 (3) 22 2 (5) 2 3 1 4 2 S 2 8 2 10 2 12 2 (11) 2 2 1 6 17 1 12 1 4 15 1 8 110 11 1 19 16 2 (7) 2 (9) 2 1 1 СС S 11 2 (3) 22 2 (5) 2 3 1 4 2 S 2 8 2 10 2 12 2 (11) 2 2 1 12 15 1 8 110 11 1 16 2 (7) 2 (9)", null, "Intersection between a Solid and a Plane. S 1(1) 2 3 2 2 2 5 2 3 1 4 2 S 2 2 1 5 1 6 2 1 1 6 14 1 1 3 3 3 2 3 (5) 34 3 S 3 (6)", null, "Surfaces. Ruled surfaces. A surface which can be generated by a straight line is called the ruled surface. Hence, the ruled surface represents a locus of straight lines. A surface which may be generated only by a curved line will be called the double- curved surface.", null, "Surfaces. Developable Ruled surfaces. Some curved surfaces can be developed so that they coincide completely (with all their points) with a plane without stretching or shrinking. Such surfaces will be called developable. They comprise only ruled surfaces in which adjacent rectilinear generatrices are parallel or intersect , or are tangent to sphere curve. All double-curved surfaces and the ruled surfaces which cannot be developed into a plane are called nondevelopable (or warped ) surfaces", null, "Surfaces. Ruled surfaces with one Diretrix. Conical Surface m 1 m 2 AS m M 2 M 1 S 2 l 1 A 1 S 1 mssmsm, , Con. S A conical surface is generated by a straight line passing through a fixed point and through all the points (in succession) of a curved directing line. This line is also called directrix or guide line. The fixed point S is called the vertex of a conical surface. S – vertex m – directrix ℓ – generatrix S 2 S", null, "Surfaces. Ruled surfaces with one Diretrix. Cylindrical surface A 2 m 1 A 1 ml 1 A ll 2 A 11 ll, , , Cyl. Siillm. Sm If vertex of conical surface is moved to infinity, then conical surface turns into cylindrical. A cylindrical surface is generated by a straight line which is parallel in all its positions to a given straight line and passes in succession through all points of a curved directing line. S – vertex → ∞ m – directrix ℓ — generatrix", null, "Surfaces. Ruled surfaces with one Diretrix. n , k – border linesmm, , TSl 1 m 1 B 1 m 2 l 2 B 2 A 1 A 2 l A m The surface called a surface with a cuspidal edge or torse surface is generated by a rectilinear generatrix performing continuous motion and touching a space curve at all its positions. This space curve serves as the directrix for the surface under consideration and is called the cuspidal edge. Obviously, by specifying the projections of the cuspidal edge, we can specify a surface in the drawing. Determinant of a torse surface : 1. Curve line m as a directrix. 2. Straight line ℓ as a generatrix. n k.", null, "Surfaces. Ruled surfaces with one Diretrix. Ruled surfaces in mechanical engineering Beveled Gear surface Spur Gear surface", null, "Surfaces. Ruled surfaces with two Diretrices. llnmnm. PS, , , where, l — generatrix , m , n — directrices ; — plane director. Ruled surface can be determined by two directrices ( guide lines ) and plane director ( plane of parallelism ). Determinant of a surface :", null, "Surfaces. Ruled surfaces with two Diretrices. CYLINDROID –– both directrices are curves. CONOID –– one directrix is a curve while the other is a straight line. WARPED PLANE –– both directrices are straight lines.", null, "Surfaces. Ruled surfaces with two Diretrices. σ 2 – plane director. l 1 l ‘ 1 m 1 M 1 l 2 M 2 m 2 l ‘ 2 σ 2 n 2 n 1 CYLINDROID The surface called a cylindroid is generated by a moving straight line ( generatrix ) which in all its positions remains parallel to a given plane (called «the plane director “ plane of parallelism ” ) and intersects two curved guide lines (two directrices). If the directrices are plane curves, then, of course, they must lie in different planes.", null, "Surfaces. Ruled surfaces with two Diretrices. σ 2 — plane director. х m 2 n 1 m 1 σ 2 n 2 M 2 M 1 CONOID The surface called a conoid is generated by a moving straight line ( generatrix ) which all the time remains parallel to a given plane (called the plane director or “ plane of parallelism ” ) and intersects two directrices one of which is a curve , the other being a straight line. If the curve is a plane one, then it must not lie in the same plane with the second directrix which is a straight line.", null, "Surfaces. Ruled surfaces with two Diretrices. σ 1 — plane director. σ х σ m 2 M 1 n 2 n 1 m 1 M 2 The surface of a warped plane ( hyperbolic paraboloid ) is determined by a plane director and two noncoplanar (skew) straight-line directrices. A straight-line generatrix moving along the directrices (and remaining parallel to the plane director) describes the surface of a hyperbolic paraboloid. It can also be obtained by planar parallel motion of one parabola as generatrix along the other parabola as directrix. WARPED PLANE (hyperbolic paraboloid)", null, "Surfaces. Ruled surfaces with two Diretrices. WARPED PLANE –– both directrices are straight lines. x m 2 n 1 m 1 m, n – directrices , П 1 – plane director. l 1 2 l 1 1 M 2 M", null, "Surfaces. Ruled surfaces with two Diretrices. Right Helicoid — ruled surface with a plane director perpendicular to one directrix — straight line (axis of helicoid) while the other directrix is a helix. Straight line directrix and axis of helix base cylinder are collinear (coinciding) lines. Thus Right helicoid surface is conoid. Algorithm: how to construct Right Helicoid. Source data : – directrix – straight line plane director ( P 1 ), m – other directrix – helix , ℓ — generatrix. • Construct two views of helix – sinusoid and circle. • Construct a few intermediate locations of generatrix ℓ which intersect both directrices and remain parallel to the plane director P 1. • Multiple locations of generatrix form the helicoid surface . Point A belongs to the helicoidal surface. Missing projection can be found by constructing proper projections of generatrix. i 1 i 2 m 2 A 2 l 2 m 1 A 1 l", null, "Surfaces. Ruled surfaces with two Diretrices. Skew Helicoid – ruled surface where generatrix ℓ moves along two directrices (one is m – helix , the other i – its axis ) and remains parallel to the base cone of revolution ( constant angle between axis and generatrix of a cone ). Curve of intersection between skew helicoid surface and a plane, perpendicular to its axis is called Archimedean spiral. Point A belongs to the skew helicoidal surface. Missing projection can be found by constructing proper projections of generatrix. i 1 i 2 m 2 A 2 l 2 m 1 A 1 l", null, "Surfaces. Ruled surfaces with two Diretrices. Right Helicoid Skew Helicoid", null, "Surfaces. Ruled surfaces with three Diretrices. 1 – general view (3 directrices – curved lines ) 2 – double-skew cylindroid (2 directrices – curved lines ) 3 – double-skew conoid (2 directrices – straight lines ) 4 – one-sheet hyperboloid (3 directrices – straight lines ) Generally Ruled Surface unambiguously represented by moving generatrix along three directrices.", null, "Surfaces. Ruled surfaces with three Diretrices. One-sheet Hyperboloid — can be formed by moving straight line generatrix along three skew lines (straight line directrices) which are not parallel to any plane. Generatrices can be directrices and vice versa. Outline — hyperbola. HYPERBOL", null, "Surfaces. Ruled surfaces with three Diretrices. Skew Wedge surface ( type of double-skew cylindroid surface ) — two directrices are smooth curves and one is a straight line. All directrices are lying in parallel planes. On construction Chords are divided proportionally. А 1 1 1 D 1 В 1 С 1 1 1 А 2 1 2 D 2 В 2 С 2 1 2 n 1 2 n 2 2 n 3 1 n 2 1 n", null, "Channel Surfaces.", null, "Channel Surfaces. O n L n. L 1 m", null, "Channel Surfaces. A 1 A m 1 O 1 lm 2 l 1 O 2 A 2 O m", null, "Channel Surfaces. 0 r=f()h 1=r h=F()", null, "Channel Surfaces. Screw Torus", null, "Surfaces. Positional problems. Intersection between a Line and a Surface. 1 1 1 21 2 2 2 T 11 1 =2 12 2 1 2 S 1 S 2 a 1 = B 1 K 1 B 2 T 2 S 1 K 2 a 2 S 2 3 Intersection between 5 — faced prismatic horizontally perpendicular surface and a straight line. Segment of a frontally perpendicular cylindrical surface intersects with a straight line Horizontally projecting straight line intersects: (а) — pyramidal surface ( b ) – conical surface", null, "Surfaces. Positional problems. Intersection between a Line and a Surface. M 1 M 2 A 1 b 2 A 2 R 1 K 1 B 1 T 11 1 4 23 2 2 1 4 1 N 1 3 1 N 2 K 21 2 2 2 T 2 R 2 B 2 x 5 b 1 Intersection between АВ -line and cylindrical surface. K – piercing point.", null, "Intersection between АВ -line and surface of revolution. Surfaces. Positional problems. Intersection between a Line and a Surface. 6 a 1 1 1 a 2 K ‘ 1 1 ‘ 2 b 1 b 2 5 ‘ 1 2 1 3 1 4 1 5 1 4 ‘ 1 3 ‘ 1 2 ‘ 1 1 ‘ 1 K ‘ 2 K 1 2 ‘ 23 ‘ 24 ‘ 25 ‘", null, "Surfaces. Positional problems. Intersection between a Line and a Surface. 81 2 N 2 1 1 x M 1 N 1 O 1 M 2 2 12 2 x’ 2 ‘ 1 M ‘ 11 ‘ 1 r O 2 N ‘ 1 7 N 114 113 2 11 2 9 1 l 1 5 17 1 1 1 M 1 M 2 3 25 2 11 1 7 29 2 l 2 1 2 6 1 2 14 18 110 112 2 12 114 2 8 2 6 24 2 2 2 1 3 1 N 2 10 2 O ‘ 1 Projecting surfaces Intersection between a sphere and oblique line", null, "Surfaces. Positional problems. Intersection between a Line and a Surface. 10 a 2 k 2 a 1 k 1 h 2 h 1 b 1 b 2 m 2 n 2 m 1 n 1 h", null, "Surfaces. Positional problems. Intersection between a Line and a Surface. n 1 12 m 1 1 1 a 1 2 1 k 1 b 1 h 1 S 1 h 2 b 2 k 2 a 2 n 2 m 2 S", null, "Surfaces. Positional problems. Intersection between surfaces. Positional problems determine relative position and mutual belongings of objects and can be solved on their orthographical views. Basic positional problem: construct curve of intersection between two surfaces. If 1– 2 surfaces are perpendicular to principal planes (projecting surfaces) , then positional problems can be solved without additional cutting surfaces.", null, "Intersection between Surfaces A 2 A 1 B 2 C 1 B 1 A 3 =C 3 B 3 S 2 S 3 S 11 2 3 2 (5) 2 (7) 2(8) 2 1 3 3 35 3 8 3 (7) 3 (6) 3 4 3 ( 2 ) 3 4 1 1 1 3 1 2 16 18 1 5 17 1 2 24 2 (6) 2 Surface perpendicular to a principal plane", null, "Intersection between Surfaces. Y 341 2 ( 2 ‘ ) 2 (3 ‘ ) 24 2 (5 ‘ ) 2 (4 ‘ ) 2 3 22 2(1 ‘ ) 2 5 2 1 1 4 1 (2) 1 (3) 1(3 ‘ ) 14 ‘ 1 5 1 (2 ‘ ) 15 ‘ 1 1 ‘ 1 Y 1 ‘Y 6 6 2 6 16 ‘ 1 ( 6 ‘ ) 2 Y 6 Y 1 ‘Y 34 Both Surfaces are perpendicular to principal planes Y 1 1 3 3 3 2 34 3 5 3 6 3 1 ‘ 33 ‘ 32 ‘ 36 ‘ 35 ‘ 3 4 ‘ 3 Circle Ellipse. Straight Line", null, "МТК Intersection between Surfaces. Polyhedrons. Both Surfaces are perpendicular to principal planes", null, "МТК Intersection between Surfaces. f 2 h 16 2 B 1 2 2 3 2 4 25 27 2 f 1 h", null, "Intersection between Surfaces. B 2 A 2 C 1 B 1 A 1 1 1 2 1 3 1 4 15 16 1 S 1 4 23 2(6) 21 2 ( 5 ) 22 2 S ‘ 2 П ‘ 2 6 ‘ 2 4 ‘ 2 5 ‘ 2 3 ‘ 2 1 ‘ 2 2 ‘ 2 K ‘ 2 L ‘ 2 M ‘ 2 N ‘ 2 O ‘ 2 P ‘ 2 B ‘ 2 A ‘ 2 =C ‘ 2 K 1 L 1 N 1 M 1 P 1 O 1 L 2 N 2 P 2 O 2 M 2 K 2 h 1 f 2 f", null, "Intersection between Surfaces. Cutting Surface Method. To construct curve of intersection between surfaces we can use Cutting Surface method. As cutting surfaces we can use Cutting Planes or Cutting Spheres. Projection of a Curve of intersection between surfaces lies in an area of intersection between projections of outlines of surfaces. To construct curve of intersection it’s desirable to indicate characteristic points : — on the outline ( upper , lower , left , right ). Ф 2 n. Ф 1 m MN", null, "Intersection between Surfaces. Cutting Plane Method. S 2 ( 2 ‘ ) 2 (3 ‘ ) 24 2 (5 ‘ ) 2 (4 ‘ ) 2 3 22 25 2 (4 ‘ ) 3 (4) 3 6 36 ‘ 3 (3) 3 2 32 ‘ 3 1 11 ‘ 1 2 1 3 13 ‘ 12 ‘ 1 4 14 ‘ 1 6 16 ‘ 1 S 3 S 1 Algorithm Determine area where curve of intersection can exist. 1. Determine which surfaces of a model are perpendicular to principal planes. 2. Define base points for curve of intersection. 3. Assign type of cutting surfaces. 4. Define maximum number of cutting surfaces and their range. 5. Construct points of a curve of intersection. 6. Define base points where visibility of a curve changes. 7. Join all obtained points by a curve line in a proper sequence. (3 ‘ ) 3 1 3 hyperbolaellipse circle 6 2 (6 ‘ ) 2 5 35 ‘ 3 1 2 (1 ‘ ) 2 5 ‘ 1 5 1 Y 3 Y", null, "Intersection between Surfaces.", null, "Intersection between Surfaces. Cutting Plane Method. R C 1 R C 2 R К 1 R К 2 R", null, "Intersection between Surfaces. Τ 1 Γ 2 2 11 1 4 23 2 Σ", null, "Intersection between Surfaces. Cutting Plane Method. а 16 1 8 1 4 14 2 9 2 3 1 Σ 1 Δ 1 = f 17 25 2 2 2 b 2 h 2 1 2 f 2 5 16 1 A 1 7 1 h ‘ 2 h » 2 а 2 B 1 D 1 S 1 2 1 h ‘ 2 b 13 28 1 A 2 B 2 C 2 1 1 Common symmetry plane S", null, "Intersection between Surfaces. Cutting Sphere Method. Coaxial surfaces of revolution (i. e. surfaces with a common axis) intersect along circles.", null, "Intersection between Surfaces. Cutting Sphere Method. Cutting surfaces — spheres Theorem: Two coaxial surfaces of revolution intersect at circles which lie in planes, perpendicular to the axis of revolution, and centers of the circles lie on that axis. If center of a sphere lies on the axes of revolution, then sphere intersects surface of revolution at circles. Application of the method : 1. Both surfaces are surfaces of revolution ; 2. Axes of these surfaces are intersecting lines ; 3. A plane, formed by intersecting axes, should be parallel to a principal plane. а. c. b. d. A B A =", null, "Intersection between Surfaces. Cutting Sphere Method. Algorithm how to construct curves of intersection between two surfaces of revolution (cones). 1. Determine point of intersection between axes of revolution as a center of cutting spheres (О 2 point ). 2. Find projections of base points for curves of intersection. 3. Define band of radii R for cutting spheres: R MAX and R MIN. R MAX – distance from the center of a sphere to the outermost point , which belongs to the curve of intersection. R MIN – radius of a sphere, which is tangent to the one surface of revolution and intersects the other one. 4. Construct sphere of a radius R: R MIN < R MAX and with a center О 2. The sphere intersect each surface at circles. These circles project on П 2 plane as straight lines. Point of intersection between these lines belongs to the curve of intersection between two surfaces. 5. Construct multiple of cutting spheres for different R. Join gained points and obtain curve of intersection between two surfaces. Cutting surfaces – concentric spheres O 2 R min R max", null, "Intersection between Surfaces. Cutting Sphere Method. Cutting surfaces – eccentric spheres Algorithm how to construct curves of intersection between two surfaces of revolution (a cone and a torus). Application of the method: both surfaces of revolution have a common plane of symmetry. 1. Determine base points 1 2 , 2 2. 2. Construct plane ∑ 2 ( sigma ), which includes axes of a torus t 2. It intersects the torus at a circle of R-radius (radius of tube). 3. On the axis of a cone find center of a sphere, which intersects torus at that circle — point О 2. 4. Draw a sphere of R’-radius and О 2 – center. 5. Find a circle of intersection between a sphere and a cone. (on the plane П 2 projects to a line ). 6. Find points of intersection between two circles – points 3(4). 7. Through the axis t 2 draw new cutting plane ∑ ‘ 2 and repeat constructions. 8. Join points 1, 3, 5, 2 by a smooth curve. R’’O ‘ 2 O 2 R ’ t 2 3 21 2 5 2 2 2Σ 2 Σ ‘ 2=(4) 2 =(6) 2 R", null, "Intersection between Surfaces. Cutting Sphere Method. Curve of intersection. Concentric spheres Eccentric spheres", null, "Intersection between Surfaces. Particular case. Monge’s theorem. In general two second-order surfaces of revolution intersect at a four-degree curve. But in some cases the curve of intersection decomposes into two planar second-degree curves. It happens when both intersecting surfaces of revolution (a cylinder and a cone, two cones, an ellipsoid and a cone, etc. ) are circumscribed about a common sphere.", null, "Intersection between Surfaces. Cutting sphere method. Particular case. Monge’s theorem. Two second-order surfaces circumscribed about a third second-order surface intersect each other along two second-order curves.", null, "Intersection between Surfaces. Cutting sphere method. Particular case. Monge’s theorem. Two second-order surfaces circumscribed about a third second-order surface intersect each other along two second-order curves.", null, "Intersection between Surfaces. Cutting sphere method. Particular case. Monge’s theorem. Curve of intersection — elliptical arcs", null, "Developments. Properties of developments. 1. Each point on the development corresponds to a single point of a surface. 2. Straight lines of a surface remain straight on a development. 3. Straight line segments preserve their lengths. 4. An angle formed by lines on a surface remains equal to an angle between the corresponding lines on the development. 5. The area of a closed domain on a surface retains its magnitude within the corresponding closed domain on the development. 6. Shortest distance (beeline) on a surface develops into a straight line on a development (developed distance). 7. Parallel lines develop to parallel lines. Some curved surfaces can be developed so that they coincide completely (with all their points) with a plane without stretching or shrinking. Obtained planar figure is called development of a surface while surfaces are called developable. They comprise polyhedrons and ruled surfaces where adjacent rectilinear generatrices are parallel or intersect, or are tangent to sphere curve. All double-curved surfaces and the ruled surfaces which cannot be developed into a plane are called nondevelopable (or warped) surfaces.", null, "Methods of exact development. 1. Triangulation method. 2. Radial-line method. 3. Stretch-out-line (right section) method. Developments. Types of developments: 1. Exact. 2. Approximate. 3. Conventional. Exact developments. Developments of polyhedrons, right circular cylinder and cone can be constructed theoretically exactly. Cylinder develops to a rectangle, cone – to a circular sector. D D HH π D L L", null, "Developments. Triangulation method. Development of a pyramid. Lateral faces of a pyramid are triangles. To find true sizes of triangles you can use transformation of a drawing or find true lengths of corresponding edges. Because base of a pyramid is parallel to П 1 -plane, it’s sufficient to find true lengths of lateral edges AS , BS , CS. After that construct, for example, В CS-face, and supplement construction adding ASB and А SC faces. S 2 C 2 B 2 A 2 B ‘ 2 C ‘ 2 A ‘ 2 S 1 C 1 B 1 A 1 A ‘ 1 B ‘ 1 C ‘ 1 S 0 C 0 B 0 A 0 B 0 R SA R SB R S С R AB R С B R BCR BA R", null, "Developments. Radial Line method. C 0 A 1 A 2 C 0 A 0 B 0 D 1 C 1 F 1 B 1 C 2 B 2 E 1 F 0 F 2 D 2 E 0 D 0 This method is used when a base of a figure is parallel to the one principal plane of projection and generatrices are parallel to the other principal plane. Find true size of А 2 D 2 F 0 C 0 –face revolving it around frontal А 2 D 2. Point F 2 moves perpendicularly to А 2 D 2 to location F 0 , which can be found by protracting А 1 В 1 – line segment (the true size of АВ ) from the point А 2. From the point B 2 draw perpendicular to А 2 D 2 and find B 0 , protracting it from the point C 0 by В 1 С 1 -line segment, etc. Thus we obtained row of points А 0 А 2 , В 0 , С 0 . . . , which define a margin for developed bottom base of a figure.", null, "Developments. Stretch out Line method. This method is used when generatrices are parallel to a principal plane of projection. Cut ABCDEF-prism by the plane , perpendicular to the lateral edges of prism. Construct section of the prism by the plane in question – 123. Define true lengths of sides of 123 -triangle. On a free space of a drawing draw horizontal straight line a. From the arbitrary point 1 0 on this straight line protract line segments [1 0 2 0 ], [2 0 3 0 ], [3 0 1 0 ], equal to the sides of 123 -triangle. Through the points 1 0 , 2 0 , 3 0 , 1 0 draw straight lines , perpendicular to the a -line , and protract from the points 1 0 , 2 0 , 3 0 , 1 0 line segments , equal to lengths of corresponding lateral edges ([1 2 A 2 ], [1 2 D 2 ], [2 2 B 2 ], [2 2 E 2 ], [3 2 С 2 ], [3 2 F 2 ]), Then join obtained points A 0 B 0 C 0 A 0 and D 0 E 0 F 0 D 0 line segments. Obtained planar figure A 0 B 0 C 0 A 0 D 0 E 0 F 0 D 0 is a development of a lateral surface of the prism. To construct full development of a prism supplement obtained development by true figures of top and bottom bases. Construct development of ABCDEF — prism. A 1 C 1 D 1 F 1 A 2 B 21 2 2 2 A ‘ 2 C ‘ 2 B ‘ 2 A 0 B 0 C 0 A 01 1 D 0 F 0 E 0 D 0 2 01 0 3 0 1 0 D ‘ 2 E ‘ 2 2 23 21 2 E 2 F 2 D 2 D 0 A 0 B 1 E 1 C 2 3 2Ƴ 2 1 3 1 F ‘", null, "Development of a conical surface. Triangulation method. 1. Substitute (approximate) the conical surface by polyhedral pyramidal surface. 2. Define true lengths of lateral edges (method of revolution is used). 3. Construct true sizes of faces. 4. Obtained points of cone base join by a smooth curve. Size and number of straight line segments which approximate a curve directrix of a cone depend on curvature of a curve and size of a cone. Approximate developments. 1 1 S 2 S 1 2 1 3 1 5 16 17 1 8 11 ‘ 2 1 ‘ 1 6 ‘ 1 5 ‘ 1 7 ‘ 1 4 ‘ 1 8 ‘ 1 3 ‘ 1 2 ‘ 16 ‘ 2 5 ‘ 2 7 ‘ 2 4 ‘ 2 8 ‘ 2 3 ‘ 2 2 ‘ 23 ‘ 2 4 ‘ 2 5 ‘ 26 ‘ 27 ‘ 28 ‘ 21 ‘ 2 r =| 1 1 2 1 | R =|S 2 2 2 |", null, "Conventional developments. Development of Toroidal surface— Zone Method Each parallel of a surface is spaced an equal distance, m , apart along the surface. Cones are passed through the surface so that they pass through two parallels at the outer surface of the toroidal. The largest cone with element R 1 is found by extending it through where the equator and the next parallel intersect on the surface in the front view until R 1 intersects the extended centerline of the toroidal surface. Elements R 2 and R 3 are found by repeating this process. The development is begun by laying out the largest zone, using R 1 as the radius, on the arc that represents the base of an imaginary cone. The breadth of the zone is found by laying off distance m from the front view to the development and drawing the upper portion of the zone with a radius equal to R 1 -m using the same center. The next zone is drawn using radius R 2 with its center located on a line through the center of arc R 1. The last zone will appear as a circle with R 3 as its radius. The lengths of the arcs can be established by dividing the top view with vertical cutting planes that radiate through the poles — meridians. Preliminary nondevelopable surface is approximated by segments of developable surfaces. Then construct development of every segment, and the whole totality of developments gives conventional development of nondevelopable surface. b c", null, "" ]
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https://www.arxiv-vanity.com/papers/hep-ph/0611158/
[ "# Connection between the Sivers function and the anomalous magnetic moment\n\nZhun Lu    Ivan Schmidt Departamento de Física, Universidad Técnica Federico Santa María, Casilla 110-V, Valparaíso, Chile\nand Center of Subatomic Physics, Valparaíso, Chile\n###### Abstract\n\nThe same light-front wave functions of the proton are involved in both the anomalous magnetic moment of the nucleon and the Sivers function. Using the diquark model, we derive a simple relation between the anomalous magnetic moment and the Sivers function, which should hold in general with good approximation. This relation can be used to provide constraints on the Sivers single spin asymmetries from the data on anomalous magnetic moments. Moreover, the relation can be viewed as a direct connection between the quark orbital angular momentum and the Sivers function.\n\n###### pacs:\n13.40.Em, 13.60.-r, 13.88.+e\npreprint: USM-TH-198\n\nThe quark orbital angular momentum seh74 (or quark transverse motion) plays an important role for understanding the spin and quark structure of the nucleon, since as shown by many studies emc88 ; jaf90 ; ji97a ; hag98 ; ma98 ; har99 , it is the missing block of the total nucleon spin. Also many interesting phenomena or observables require the presence of quark orbital motion, among which the Sivers function sivers has attracted a lot of interest, since it is an essential piece in our understanding of the single spin asymmetries (SSA) observed in semi-inclusive deeply inelastic scattering (SIDIS). These SSAs have been measured recently by both the HERMES Airapetian:2004tw ; hermes05 and COMPASS compass ; compass06 Collaborations. Denoted as , the Sivers function describes the unpolarized distribution of the quark inside a transversely polarized nucleon, which comes from a correlation of the nucleon transverse spin and the quark transverse momentum. Although this is a -odd type correlation, it has been found that final state interaction bhs (FSI) between the struck quark and the spectator system can produce the necessary phase for a non-zero Sivers function, and its QCD definition collins02 ; jy ; belitsky ; bmp03 has just been worked out. Besides the single spin asymmetry, another important feature of the Sivers function is that it encodes the parton’s orbital angular momentum () inside the nucleon. This comes from the fact that the Sivers function requires the nucleon helicity to be flipped from the initial to the final state, while the quark helicity remains unchanged. A convenient tool to study this kind of single spin asymmetry (or the corresponding Sivers function) is the light-front formalism bro98 , which can express the Sivers function as the overlap integration of light-front wavefunctions differing by  bhs ; bg06 . The same kind of overlap integration bl80 ; bhms ; bdh of light-front wavefunctions (with in the initial and final states) also appears in the anomalous magnetic moment of the nucleon, which apparently encodes the quark orbital angular momentum bhms . Therefore, it is interesting to find relations between the Sivers function and the anomalous magnetic moment of the nucleon, which is the main goal of this work. With such a relation one can constrain the Sivers function and the related asymmetries from data on nucleon anomalous magnetic moments, and viceversa. Also, the relation can be viewed as a direct connection between the quark orbital angular momentum distribution and the Sivers function.\n\nThe proton state can be expanded in a series of Fock states with coefficients , which are the light-front wavefunctions of the proton:\n\n Ψp(P+,P−,0⊥)=∑nψn(xi,k⊥,λi)|n,xiP+,k⊥,λi⟩. (1)\n\nHere is the light-front momentum fraction of the quark, denotes the helicity, and its transverse momentum. The wavefunctions are Lorentz-invariant functions of the kinematics of the constituents of nucleon, and , with and .\n\nAs pointed out before, the Sivers function requires that the nucleon wavefunctions in the initial and final state differ by one unit of orbital angular momentum, and final state interactions play a crucial role. It describes the interference of two amplitudes which have different initial proton spin but couple to the same final state: Im. This can be realized by a gluon exchange between the struck quark and the spectator system. There have been already attempts jmy02 ; bg06 , using the proton light-front wavefunctions, to find a formula to calculate the Sivers functions. In those papers the final state interaction phase needed for Sivers functions has been incorporated into the wavefunctions. Another possibility is to express the Sivers function as the product of wavefunctions and the final state interactions term:\n\n f⊥1T ∝ ∑nψ↑∗n(xi,k⊥,λi)G(xi,x′i,k⊥,i,k′⊥,i)ψ↓n(x′i,k′⊥,λ′i), (2)\n\nwhere is the final state interaction term, and the light-front wavefunctions in this equation are the usual ones which do not contain the final state interactions phase.\n\nIn the scalar diquark model, the calculation gets simplified and some instructive result can be derived. In this model the kinematics is determined by the momentum fraction and the transverse momentum of the struck quark , since those of the spectator diquark can be related to by , . The advantage of the diquark model is that it is simple, and it works quite well in describing other properties of the proton, such as the helicity distribution ma96 , the transversity ma00 and the form factors ma02 of the proton. The light-front wavefunctions of this two-body Fock state have the simple forms bl80 ; bhs :\n\n ⎧⎪ ⎪ ⎪⎨⎪ ⎪ ⎪⎩ψ↑+12(x,k⊥)=(M+mx)φ(x,k⊥) ,ψ↑−12(x,k⊥)=−(+k1+ik2)xφ(x,k⊥) . (3)\n\nfor and\n\n ⎧⎪ ⎪ ⎪⎨⎪ ⎪ ⎪⎩ψ↓+12(x,k⊥)=(+k1−ik2)xφ(x,k⊥) ,ψ↓−12(x,k⊥)=(M+mx)φ(x,k⊥) . (4)\n\nfor . Here\n\n φ(x,k⊥)=g/√1−xM2−(k⊥+m2)/x−(k⊥+M2d)/(1−x), (5)\n\nand , and are the masses of the proton, the quark and the diquark, respectively.\n\nAccording to Fig. 1, a formula to calculate the Sivers function can be given by the overlap integration of the proton light-front wavefunction:\n\n kL2Mf⊥q1T(x,k⊥) = i∑n∫d2k′⊥dx′16π3ψ↑⋆n(x,k⊥,λi)G(x,x′,k⊥,k′⊥)ψ↓n(x′,k′⊥,λ′i), (6)\n\nwhere , and is the kernel which contain the final state interaction. In the scalar diquark model, the kernel has a simple form bhs ; bh04 :\n\n G(x,x′,k⊥,k′⊥)=iCFαsδ(x−x′)2π((k⊥−k′⊥)2+λ2g), (7)\n\nwhich can be calculated in the eikonal approximation. Final state interactions arise from gluonic exchanges between the struck quark and the spectator system, which are necessary in order to insure gauge invariance of the -dependent distributions. In our case we use a one-gluon exchange (with momentum ) approximation, which is commonly used in most model calculations bhs ; jy ; gg02 ; bbh03 ; yuan ; bsy04 ; lm04a . Here is the mass of the gluon, which is needed in general to regularize the infrared divergence in the integration, while in our case it can be set to . The sign in Eq. (7) is according to the Trento conventions tc . With Eqs. (3), (4) and (6), one can obtain the result for Sivers function that has given in jy ; bbh . An expression similar to (6) has been given in Ref. bh04 , where the Sivers function (or the SSA) is expressed as a product of a generalized parton distribution (GPD) and a FSI term in impact parameter space, while our result is expressed in momentum space, which can be connected with the result in Ref. bh04 by Fourier transformation.", null, "Figure 1: Diagram to calculate the Sivers function. The arrows ↑/↓ show the polarizations of the nucleon.\n\nThe same set of wavefunctions (spin-flipped in the initial and final states) appearing in (2) also appears in the expression of the Pauli form factor or anomalous magnetic momentum of the proton. In light-front formalism the Pauli form factor is identified from the helicity-flip vector current matrix elements of the current bhms\n\n ⟨P+q,↑∣∣∣J+(0)2P+∣∣∣P,↓⟩=−qLF2(q2)2M. (8)\n\nFurthermore, from Fig. 2, one can express the Pauli form factor in terms of the light-front wavefunctions as\n\n −qLF2(q2)2M = ∑n∫d2k⊥dx16π3∑jejψ↑⋆n(xi,k′⊥,λi)ψ↓n(xi,k⊥,λi), (9)\n\nwhere\n\n k′⊥=k⊥+(1−x)q⊥, (10)\n\nhere is momentum of the virtual photon, the momentum transfer between the initial and final nucleon.\n\nThe anomalous magnetic moment is defined from the Pauli form factor in the limit of : . After some algebra calculation, the anomalous magnetic moment can be expressed in terms of a local matrix element at zero momentum transfer:\n\n κM = −∑jej∑n∫d2k⊥dx16π3∑i≠jψ↑⋆n(xi,k⊥,λi)xi(∂∂k1+i∂∂k2)ψ↓n(xi,k⊥,λi). (11)\n\nFrom (2) and (11) we see that the same set of light-front amplitudes, with the orbital angular momenta differing by between the initial and final state, appears in the calculation of the Sivers function and the anomalous magnetic moment. Recently a study burkardt06 has shown how the non-vanishing anomalous magnetic moment constrains the quark orbital angular momentum.\n\nIt is clear from the previous expressions that the relation we are seeking between the Sivers function and the anomalous magnetic moment will necessarily have to be approximate. This can be also seen from the fact that the Sivers function, although leading twist, has () scale evolution dependence, which is not the case for the anomalous magnetic moment.\n\nIn the scalar diquark model, the calculation of the anomalous magnetic moment at the quark level has been given in bhms as\n\n κM = ∑n=±1/2∫d2k⊥dx16π3ψ↑⋆n(x,k⊥)(1−x)(∂∂k1+i∂∂k2)ψ↓n(x,k⊥). (12)\n\nUsing the wavefunctions given in (3) and (4), one gets the result\n\n κ = Mg28π2∫10dx(1−x)2(m+xM)(1−x)m2+xm2d−x(1−x)M2 (13) = ∫10dx(1−x)AB,\n\nwhere we have defined:\n\n A = g28π2(1−x)(xM+m); (14) B = (1−x)m2+xM2d−x(1−x)M2. (15)", null, "Figure 2: Diagram to calculate the Pauli form factor (or anomalous magnetic moment). The arrows ↑/↓ show the polarizations of the nucleon.\n\nOne can calculate the lowest -moment of defined as ()\n\n f⊥,q1T(x)=∫d2k⊥f⊥,q1T(x,k2⊥). (16)\n\nAccording to Eq. (6), which gives , and the wavefunctions given in (3) and (4), we directly get\n\n f⊥,q1T(x)=−π2CFαsA12B. (17)\n\nDefining\n\n κq(x)=(1−x)AB (18)\n\nwhich satisfies the normalization condition\n\n ∫10dxκq(x)=κq, (19)\n\nwe arrive at a simply relation between and in the quark-diquark model:\n\n f⊥,q1T(x)=−π2CFαs12(1−x)κq(x). (20)\n\nThe term in this equation is consistent with the result in Ref. bg06 . From this we can get:\n\n ∫10dx(1−x)f⊥,q1T(x) = −∫10dxπ2CFαs12κq(x) (21) = −π2CFαs12κq.\n\nAn approximate relation can also be obtained from the above equation, as\n\n ⟨1−x⟩F⊥,q1T = −π2CFαs12κq, (22)\n\nwhere is the first -moment of defined as , and is the average value of in the proton. Therefore we obtain a interesting relation that provides a constraint on the Sivers function from anomalous magnetic moment data. It suggests that the moment of the Sivers function is proportional to the anomalous magnetic moment contributed by the same quark. Although this relation as been obtained in the diquark model, it should hold in general with good approximation. This can be seen from the general relations for the anomalous magnetic moment and for the Sivers function in terms of light-front wavefunctions, as given by Eqs. (6) and (11), respectively. Since both quantities calculated from the same wavefunctions, and these are the ones that vary more strongly in the range of values covered by the integrals, it is appropriate to apply the mean value theorem, and therefore get the same type of relation as the one we have found here. In this sense our approximate result can be considered more general than the diquark model, and thus be applied to each quark individually.\n\nA more fundamental object than the anomalous magnetic moment is the spin-flipped generalized parton distribution (GPD)  dm ; ji97a ; ji97b ; dvcs . Actually the function is at the forward limit ( and ):\n\n κq(x)=Eq(x,0,0). (23)\n\nThus in the scalar diquark model, is proportional to . According to Ji’s sum rule ji97a ( is the total angular momentum carried by quark flavor ):\n\n ∫dxx(Hq(x,ξ,t)+Eq(x,ξ,t))=12Jq(t), (24)\n\nwhich also holds in the forward limit. We see that the Sivers function is related to the angular momentum of the parton inside the nucleon, and therefore it is in fact sensitive to the orbital angular momentum of the quark. There is then the possibility to get information of the quark orbital angular momentum from the Sivers functions.\n\nAlthough the relation given in (20) and (21) is a simple result based on the approximation of the scalar diquark model, we can still apply the relation to given some prediction on the Sivers asymmetry of the meson production in SIDIS processes, such as the ratio of the asymmetries between different final mesons, since in this case the model dependence is reduced. The Sivers asymmetry is proportional to\n\n ASivUT∝⟨∑ae2af⊥a1TDa1⟩⟨∑ae2afa1Da1⟩, (25)\n\nwhich can be extracted by the weighting function , here and denote respectively the azimuthal angles of the produced hadron and of the nucleon spin polarization, with respect to the lepton scattering plane, is the unpolarized fragmentation function. We will focus on the large regime that the valence quark contribution dominates, and the disfavored fragmentation function can be ignored.\n\nSince the Sivers function and the anomalous magnetic moment “share” the same set of the proton wavefunctions, as shown in Figs. 1 and 2, one can start from the data of the anomalous magnetic moment to provide constraints on the proton wavefunctions, and then on the Sivers function. Similar methods have been used in Ref. bg06 , where a small Sivers asymmetry on a deuteron target has been predicted, and Ref. burkardt02 , where the sign of the Sivers asymmetries for different hadron targets combining different fragmenting hadrons has been classified.\n\nAs figured out in Ref. bg06 , the quark contribution dominates over the gluon contribution in the case of Sivers functions Jacques , which is also the result of an argument based on large considerations pobylista . There are also phenomenological supports of this conclusion from the SIDIS experiment from COMPASS of CERN compass , as pointed out in Ref. bg06 , and the analysis on hadron production of given in Ref. ans06 . Therefore in this work we only consider the quark contribution to the Sivers functions and the corresponding asymmetry.\n\nOne can constrain the proton wavefunctions by normalizing each and quarks contributions to the anomalous moments , . Isospin symmetry implies\n\n κd/n = κu/p, (26) κu/n = κd/p. (27)\n\nIn the valence quark approximation we have:\n\n κp = (2)(2/3)κu/p+(−1/3)κd/p, (28) κn = (2)(−1/3)κu/p+(2/3)κd/p. (29)\n\nThus one has , . In the following we use and to represent and , respectively.\n\nThen we can write the ratio of the asymmetries between and at large :\n\n ASivUT(π+)ASivUT(π−)≈2e2uf⊥u1TDπ+/u1e2df⊥d1TDπ−/d1≈2e2uκue2dκd=−3.3. (30)\n\nAlso we have\n\n ASivUT(π0)ASivUT(π−) ≈ 2e2uf⊥u1TDπ0/u1+e2df⊥d1TDπ0/d1e2df⊥d1TDπ−/d1 (31) ≈ 2e2uκu+e2dκd2e2dκd=−1.15, ASivUT(K+)ASivUT(K0) ≈ 2e2uf⊥u1TDK+/u1e2df⊥d1TDK0/d1≈4e2uκue2dκd=−6.6. (32)\n\nFor the above result we used isospin symmetry for the quark fragmentation functions:\n\n Dπ+/u1 = Dπ−/d1=2Dπ0/u1=2Dπ0/d1, (33) DK+/u1 = 2DK0/d1. (34)\n\nThe results show that in the large region, the Sivers asymmetry of is 3 times larger than that of and with opposite sign, which is consistent with the recent HERMES result hermes05 where at a four times larger asymmetry of is measured; the asymmetries of and are similar in size; the asymmetry of is much larger than that of , nearly one order of magnitude; and since in valence approximation.\n\nIn summary, both the formula for calculating Sivers function and that for calculating the anomalous magnetic moment of the proton, can be expressed in terms of the same set of the light-front wavefunctions, with helicity flipped from initial state to final states. Using the overlap representations of both Sivers functions as well as the anomalous magnetic moment, we give a simple relation between these two observables, in the approximation of the scalar diquark model. This relation is applied to provide constraints on the Sivers single spin asymmetries in the valence regime from data on anomalous magnetic moments. Also, the relation can be viewed as a connection between the quark orbital angular momentum and the Sivers function.\n\nAcknowledgements. This work is supported by Fondecyt (Chile) under Project No. 3050047, and by the ”Center of Subatomic Physics” (Chile) ." ]
[ null, "https://media.arxiv-vanity.com/render-output/5663010/x1.png", null, "https://media.arxiv-vanity.com/render-output/5663010/x2.png", null ]
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http://www-groups.dcs.st-and.ac.uk/~history/Biographies/Ackermann.html
[ "Wilhelm Ackermann\n\nBorn: 29 March 1896 in Schönebeck (Kr. Altena), Germany Died: 24 December 1962 in Lüdenscheid, Germany", null, "Click the picture above\nto see a larger version\n\nWilhelm Ackermann was a mathematical logician who worked with David Hilbert in Göttingen but spent his career as a high school teacher.\n\nAckermann entered the University of Göttingen in 1914 to study mathematics, physics and philosophy. However soon after he began his studies Word War I started. Ackermann was drafted into the army in 1915 and continued to serve until 1919 when he was able to return to his studies in Göttingen. He received his doctoral degree in 1925 with a thesis Begründung des \"tertium non datur\" mittels der Hilbertschen Theorie der Widerspruchsfreiheit written under David Hilbert's supervision. It provided a proof of the consistency of arithmetic without induction. It was intended to be a consistency proof for elementary analysis although this proof contained significant errors. Richard Zach explains the background to the thesis in :-\n\nAckermann's 1924 dissertation is of particular interest since it is the first non-trivial example of what Hilbert considered to be a finitistic consistency proof. Von Neumann's paper of 1927, the only other major contribution to proof theory in the 1920s, does not entirely fit into the tradition of the Hilbert school, and we have no evidence of the extent of Hilbert's involvement in its writing. Later consistency proofs, in particular those by Gentzen and Kalmár, were written after Gödel's incomplete ness results were already well-known and their implications understood by proof theorists. Ackermann's work, on the other hand, arose entirely out of Hilbert's research project, and there is ample evidence that Hilbert was aware of the range and details of the proof.\nAfter submitting his dissertation, Ackermann went to Cambridge, England, where he spent the first half of 1925. He was awarded a Rockefeller Scholarship to support this trip and Hilbert wrote in support of his application showing his high opinion of Ackerman's work :-\nIn his thesis \"Foundation of the 'tertium non datur' using Hilbert's theory of consistency,\" Ackermann has shown in the most general case that the use of the words \"all\" and \"there is,\" of the \"tertium non datur,\" is free from contradiction. The proof uses exclusively primitive and finite inference methods. Everything is demonstrated, as it were, directly on the mathematical formalism. Ackermann has here surmounted considerable mathematical difficulties and solved a problem which is of first importance to the modern efforts directed at providing a new foundation for mathematics.\nAckermann was also the main contributor to the development of the logical system known as the epsilon calculus, originally due to Hilbert. This formalism formed the basis of Bourbaki's logic and set theory.\n\nFrom 1929 until 1948 he taught as a teacher at the Arnoldinum Gymnasium in Burgsteinfurt and in Luedenscheid. He was corresponding member of the Akademie der Wissenschaften in Göttingen, and was honorary professor at the Universität Münster.\n\nIn 1928, Ackermann observed that A(x, y, z), the z-fold iterated exponentiation of x with y, is an example of a recursive function which is not primitive recursive. A(x, y, z) was simplified to a function P(x, y) of two variables by Rozsa Peter whose initial condition was simplified by Raphael Robinson. It is the latter which occurs as Ackermann's function in today's textbooks. Also in 1928 the often reprinted book Grundzüge der Theoretischen Logik by Hilbert and Ackermann appeared.\n\nAmong Ackermann's later work are consistency proofs for set theory (1937), full arithmetic (1940) and type free logic (1952). Further there was a new axiomatization of set theory (1956), and a book Solvable cases of the decision problem (North Holland, 1954. Second edition, 1962).\n\nThe new axiomatization of set theory was presented by Ackermann in Zur Axiomatik der Mengenlehre (1956). Dana Scott writes :-\n\nA remarkably simple axiomatization of a system of set theory is presented which the reviewer feels deserves serious consideration. The system is formalized in an applied first-order calculus with identity using a binary predicate e (membership) and a singulary predicate M (being a set). It is quite essential for the consistency of the system that M is not definable in terms of E. The axiom of extensionality is assumed so that all the individuals can be considered as collections of individuals, but it is easily proved from the axioms that there are collections that are not sets and even contain non-sets. The necessity for the existence of such improper collections in the theory makes comparison with the standard systems of set theory somewhat difficult.\nRudolf Grewe, who wrote his doctoral dissertation on Ackermann's set theory in 1966, gives models for the theory in . Several other authors have studied this system and references are given in .\n\nIn 1957 Ackermann published Philosophische Bemerkungen zur mathematischen Logik und zur mathematischen Grundlagenforschung and its English translation Philosophical observations on mathematical logic and on investigations into the foundations of mathematics. This paper, written for non-experts in the subject, gives an excellent overview of how Ackermann viewed mathematical logic. John Van Heijenoort writes :-\n\nAn objection to mathematical logic is that it is not the same as the philosophical logic which forms the foundation of our thought and which alone is necessary for thinking. Ackermann remarks that the traditional modes of inference are included in mathematical logic, besides many others, like the statement logic or the logic of relations. One has the illusion of getting by with \"Aristotelian logic\" in mathematics just as long as mathematical reasonings are insufficiently analyzed.\n\nA further objection to mathematical logic is that it is incomplete, in the sense that, by Gödel's result of 1931 (the text says 1932), intuitively correct theorems cannot be proved within a given system. But intuitive thinking is not consistent; paradoxes arise on the basis of modes of inference which from the naive intuitive point of view have to be considered correct. If we restrict intuitive inference so that paradoxes are eliminated, we come to a formal system, and incompleteness is the price we have to pay for consistency.\n\nMathematical logic is of use in clarifying the concepts of \"necessity\" and \"possibility,\" the distinction between \"analytic\" and \"synthetic.\" Discussing the \"triviality\" of logic, the author presents the decision problem.\n\nMathematics is today viewed as an investigation cf structures; but the obviousness connected with the concept of natural number is independent of any axiomatically introduced structure. Ackermann presents intuitionism, which constructs a mathematics with a minimum of logic, and the Frege-Russell analysis of the number concept. He defends this analysis against certain objections (circularity, necessity of an axiom of infinity). He remarks that Brouwer reduced the number intuition to two concepts: unity and possibility of repeatedly distinguishing a unity from another. In these concepts \"there is nothing foreign to logic.\" He concludes that \"in the theory of natural numbers we have a domain which is capable of an intuitive foundation with a minimum of logic in the sense of [Brouwer] and also of attainment purely by logical definitions if we presuppose an extensive logic.\"\n...\nThe paper ends with remarks on geometry, as a science conveying knowledge of the external world. Beyond any axiomatic treatment of geometry there are intuitive geometrical impressions which force themselves upon us with compelling power. In order to obtain a systematic and complete whole we may add principles in which free constructions occur. But, although limits may be difficult to trace, there remains a nucleus of obligatory intuitive elements. Thus the author sketches a position distinct, on the one hand, from a comprehensive apriorism in the sense of Kant and, on the other hand, from a thoroughgoing conventionalism.\n\nArticle by: J J O'Connor, E F Robertson, and Walter Felscher, Tübingen\n\nClick on this link to see a list of the Glossary entries for this page\n\nList of References (6 books/articles)\n\nMathematicians born in the same country\n\nOther Web sites\n1. Encyclopedia of Philosophy (The epsilon-calulus)\n2. World History\n3. Stanford Encyclopedia of Philosophy (The epsilon-calulus)\n4. Mathematical Genealogy Project" ]
[ null, "http://www-groups.dcs.st-and.ac.uk/~history/Thumbnails/Ackermann.jpg", null ]
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https://akjournals.com/search?access=all&page=15&pageSize=10&q=%22thermodynamic+functions%22&sort=relevance
[ "# Search Results\n\n## You are looking at 141 - 150 of 215 items for :\n\n• \"thermodynamic functions\"\n• Refine by Access: All Content\nClear All\n\n## Solvent extraction of thorium from nitric acid solutions using di- N -butyl sulfoxide (DBSO) in xylene\n\nJournal of Radioanalytical and Nuclear Chemistry\nAuthors:\n,\nSyed Hasany\n, and\nAkbar Ali\n\n## Abstract\n\nThe extraction of thorium(IV) from nitric acid solutions by di-n-butyl sulfoxide (DBSO) in xylene has been investigated as a function of acid, extractant and the metal concentration. The effect of contact time and diverse ions on the extraction has been examined. Phosphate, fluoride, oxalate and perchlorate reduce the extraction to some extent. The extraction of other metal ions, especially impurities associated with thorium in ores, has been measured under optimised conditions selected for thorium extraction. Na(I), K(I), Ca(II), Sr(II), Mn(II), Fe(II), Ni(II), Zn(II), Pb(II), Al(III), Ti(IV) and Hf(IV) are not extracted. Among the stripping solutions employed for back-extraction, deionized water is found to be the best and more than 99% thorium can be back-extracted in three stages. The extracted species is supposed to be Th(NO3)4·2DBSO. The extraction is found to be almost independent of the thorium concentration in the range between 4.3·10–4–4.3·10–2M and inversely dependent upon the temperature. The values of thermodynamic functions", null, "H,", null, "G and", null, "S for extraction equilibrium have been evaluated to be –19.6±2.9 kJ·mole–1, –18.1±2.0 kJ·mole–1 and –5.0±2.9 J·mole–1·K–1, respectively.\n\nRestricted access\n\n## Molar heat capacities and standard molar enthalpy of formation of 2-amino-5-methylpyridine\n\nJournal of Thermal Analysis and Calorimetry\nAuthors:\nJ. Zhang\n,\nZ. Tan\n,\nQ. Meng\n,\nQ. Shi\n,\nB. Tong\n, and\nS. Wang\n\n## Abstract\n\nThe heat capacities (C p,m) of 2-amino-5-methylpyridine (AMP) were measured by a precision automated adiabatic calorimeter over the temperature range from 80 to 398 K. A solid-liquid phase transition was found in the range from 336 to 351 K with the peak heat capacity at 350.426 K. The melting temperature (T m), the molar enthalpy (Δfus H m 0), and the molar entropy (Δfus S m 0) of fusion were determined to be 350.431±0.018 K, 18.108 kJ mol−1 and 51.676 J K−1 mol−1, respectively. The mole fraction purity of the sample used was determined to be 0.99734 through the Van’t Hoff equation. The thermodynamic functions (H T-H 298.15 and S T-S 298.15) were calculated. The molar energy of combustion and the standard molar enthalpy of combustion were determined, ΔU c(C6H8N2,cr)= −3500.15±1.51 kJ mol−1 and Δc H m 0 (C6H8N2,cr)= −3502.64±1.51 kJ mol−1, by means of a precision oxygen-bomb combustion calorimeter at T=298.15 K. The standard molar enthalpy of formation of the crystalline compound was derived, Δr H m 0 (C6H8N2,cr)= −1.74±0.57 kJ mol−1.\n\nRestricted access\n\n## Thermodynamic investigation of several natural polyols\n\n### Part III. Heat capacities and thermodynamic properties of erythritol\n\nJournal of Thermal Analysis and Calorimetry\nAuthors:\nB. Tong\n,\nZ. Tan\n,\nJ. Zhang\n, and\nS. Wang\n\n## Abstract\n\nThe low-temperature heat capacity C p,m of erythritol (C4H10O4, CAS 149-32-6) was precisely measured in the temperature range from 80 to 410 K by means of a small sample automated adiabatic calorimeter. A solid-liquid phase transition was found at T=390.254 K from the experimental C p-T curve. The molar enthalpy and entropy of this transition were determined to be 37.92±0.19 kJ mol−1 and 97.17±0.49 J K−1 mol−1, respectively. The thermodynamic functions [H T-H 298.15] and [S T-S 298.15], were derived from the heat capacity data in the temperature range of 80 to 410 K with an interval of 5 K. The standard molar enthalpy of combustion and the standard molar enthalpy of formation of the compound have been determined: Δc H m 0(C4H10O4, cr)= −2102.90±1.56 kJ mol−1 and Δf H m 0(C4H10O4, cr)= − 900.29±0.84 kJ mol−1, by means of a precision oxygen-bomb combustion calorimeter at T=298.15 K. DSC and TG measurements were performed to study the thermostability of the compound. The results were in agreement with those obtained from heat capacity measurements.\n\nRestricted access\n\n## Rapid synthesis, kinetics and thermodynamics of binary Mn0.5Ca0.5(H2PO4)2 · H2O\n\nJournal of Thermal Analysis and Calorimetry\nAuthors:\nBanjong Boonchom\nand\nChanaiporn Danvirutai\n\n## Abstract\n\nThe binary manganese and calcium dihydrogen phosphate monohydrate Mn0.5Ca0.5(H2PO4)2 · H2O was synthesized by a rapid and simple co-precipitation method using phosphoric acid, manganese metal, and calcium carbonate at ambient temperature. Thermal transformation shows complex processes and the final decomposed product was the binary manganese calcium cyclotetraphosphate MnCaP4O12. The activation energies of some decomposed steps were calculated by Kissinger method. Activated complex theory has been applied to each step of the reactions and the thermodynamic functions are calculated. These values for transformation stages showed that they are non-spontaneous processes without the introduction of heat. The differences of physical and chemical properties of the synthesized compound and its decomposed product are compared with M(H2PO4)2 · H2O and M2P4O12 (M = Mn and Ca), which indicate the effects of the presence of Ca ions in substitution of Mn ions and confirm the formation of solid solution.\n\nRestricted access\n\n## Kinetic and thermodynamic studies of MgHPO4 · 3H2O by non-isothermal decomposition data\n\nJournal of Thermal Analysis and Calorimetry\nAuthor:\nBanjong Boonchom\n\n## Abstract\n\nThe thermal decomposition of magnesium hydrogen phosphate trihydrate MgHPO4 · 3H2O was investigated in air atmosphere using TG-DTG-DTA. MgHPO4 · 3H2O decomposes in a single step and its final decomposition product (Mg2P2O7) was obtained. The activation energies of the decomposition step of MgHPO4 · 3H2O were calculated through the isoconversional methods of the Ozawa, Kissinger–Akahira–Sunose (KAS) and Iterative equation, and the possible conversion function has been estimated through the Coats and Redfern integral equation. The activation energies calculated for the decomposition reaction by different techniques and methods were found to be consistent. The better kinetic model of the decomposition reaction for MgHPO4 · 3H2O is the F 1/3 model as a simple n-order reaction of “chemical process or mechanism no-invoking equation”. The thermodynamic functions (ΔH*, ΔG* and ΔS*) of the decomposition reaction are calculated by the activated complex theory and indicate that the process is non-spontaneous without connecting with the introduction of heat.\n\nRestricted access\n\n## Low-temperature heat capacities and thermodynamic properties of 2,2-dimethyl-1,3-propanediol\n\nJournal of Thermal Analysis and Calorimetry\nAuthors:\nB. Tong\n,\nZ. Tan\n,\nX. Lv\n,\nL. Sun\n,\nF. Xu\n,\nQ. Shi\n, and\nY. Li\n\n## Abstract\n\nThe molar heat capacities C p,m of 2,2-dimethyl-1,3-propanediol were measured in the temperature range from 78 to 410 K by means of a small sample automated adiabatic calorimeter. A solid-solid and a solid-liquid phase transitions were found at T-314.304 and 402.402 K, respectively, from the experimental C p-T curve. The molar enthalpies and entropies of these transitions were determined to be 14.78 kJ mol−1, 47.01 J K−1 mol for the solid-solid transition and 7.518 kJ mol−1, 18.68 J K−1 mol−1 for the solid-liquid transition, respectively. The dependence of heat capacity on the temperature was fitted to the following polynomial equations with least square method. In the temperature range of 80 to 310 K, C p,m/(J K−1 mol−1)=117.72+58.8022x+3.0964x 2+6.87363x 3−13.922x 4+9.8889x 5+16.195x 6; x=[(T/K)−195]/115. In the temperature range of 325 to 395 K, C p,m/(J K−1 mol−1)=290.74+22.767x−0.6247x 2−0.8716x 3−4.0159x 4−0.2878x 5+1.7244x 6; x=[(T/K)−360]/35. The thermodynamic functions H TH 298.15 and S TS 298.15, were derived from the heat capacity data in the temperature range of 80 to 410 K with an interval of 5 K. The thermostability of the compound was further tested by DSC and TG measurements. The results were in agreement with those obtained by adiabatic calorimetry.\n\nRestricted access\n\n## Low-temperature heat capacities and standard molar enthalpy of formation of the complex\n\n### Zn(Val)SO4·H2O(s) (Val=L-α-valine)\n\nJournal of Thermal Analysis and Calorimetry\nAuthors:\nY. Y. Di\n,\nZ. C. Tan\n,\nL. W. Li\n,\nS. L. Gao\n, and\nL. X. Sun\n\n## Abstract\n\nLow-temperature heat capacities of a solid complex Zn(Val)SO4·H2O(s) were measured by a precision automated adiabatic calorimeter over the temperature range between 78 and 373 K. The initial dehydration temperature of the coordination compound was determined to be, T D=327.05 K, by analysis of the heat-capacity curve. The experimental values of molar heat capacities were fitted to a polynomial equation of heat capacities (C p,m) with the reduced temperatures (x), [x=f (T)], by least square method. The polynomial fitted values of the molar heat capacities and fundamental thermodynamic functions of the complex relative to the standard reference temperature 298.15 K were given with the interval of 5 K.\n\nEnthalpies of dissolution of the [ZnSO4·7H2O(s)+Val(s)] (Δsol H m,l 0) and the Zn(Val)SO4·H2O(s) (Δsol H m,2 0) in 100.00 mL of 2 mol dm−3 HCl(aq) at T=298.15 K were determined to be, Δsol H m,l 0=(94.588±0.025) kJ mol−1 and Δsol H m,2 0=–(46.118±0.055) kJ mol−1, by means of a homemade isoperibol solution–reaction calorimeter. The standard molar enthalpy of formation of the compound was determined as: Δf H m 0 (Zn(Val)SO4·H2O(s), 298.15 K)=–(1850.97±1.92) kJ mol−1, from the enthalpies of dissolution and other auxiliary thermodynamic data through a Hess thermochemical cycle. Furthermore, the reliability of the Hess thermochemical cycle was verified by comparing UV/Vis spectra and the refractive indexes of solution A (from dissolution of the [ZnSO4·7H2O(s)+Val(s)] mixture in 2 mol dm−3 hydrochloric acid) and solution A’ (from dissolution of the complex Zn(Val)SO4·H2O(s) in 2 mol dm−3 hydrochloric acid).\n\nRestricted access\n\n## Low-temperature thermodynamics of Ln(Me2dtc)3(C12H8N2) (Me2dtc = dimethyldithiocarbamate, Ln = La, Pr, Nd, Sm)\n\nJournal of Thermal Analysis and Calorimetry\nAuthors:\nJian Wu\n,\nSan-Ping Chen\n,\nYou-Ying Di\n, and\nSheng-Li Gao\n\n## Abstract\n\nThe heat capacities of Ln(Me2dtc)3(C12H8N2) (Ln = La, Pr, Nd, Sm, Me2dtc = dimethyldithiocarbamate) have been measured by the adiabatic method within the temperature range 78–404 K. The temperature dependencies of the heat capacities, C p,m [La(Me2dtc)3(C12H8N2)] = 542.097 + 229.576 X − 27.169 X 2 + 14.596 X 3 − 7.135 X 4 (J K−1 mol−1), C p,m [Pr(Me2dtc)3(C12H8N2)] = 500.252 + 314.114 X − 17.596 X 2 − 0.131 X 3 + 16.627 X 4 (J K−1 mol−1), C p,m [Nd(Me2dtc)3(C12H8N2)] = 543.586 + 213.876 X − 68.040 X 2 + 1.173 X 3 + 2.563 X 4 (J K−1 mol−1) and C p,m [Sm(Me2dtc)3(C12H8N2)] = 528.650 + 216.408 X − 16.492 X 2 + 12.076 X 3 + 4.912 X 4 (J K−1 mol−1), were derived by the least-squares method from the experimental data. The heat capacities of Ce(Me2dtc)3(C12H8N2) and Pm(Me2dtc)3(C12H8N2) at 298.15 K were evaluated to be 617.99 and 610.09 J K−1 mol−1, respectively. Furthermore, the thermodynamic functions (entropy, enthalpy and Gibbs free energy) have been calculated using the obtained experimental heat capacity data.\n\nRestricted access\n\n## Thermodynamic investigation of room temperature ionic liquid\n\n### The heat capacity and thermodynamic functions of BMIPF6\n\nJournal of Thermal Analysis and Calorimetry\nAuthors:\nZ. Zhang\n,\nT. Cui\n,\nJ. Zhang\n,\nH. Xiong\n,\nG. Li\n,\nL. Sun\n,\nF. Xu\n,\nZ. Cao\n,\nF. Li\n, and\nJ. Zhao\n\n## Abstract\n\nThe molar heat capacities of the room temperature ionic liquid 1-butyl-3-methylimidazolium hexafluoroborate (BMIPF6) were measured by an adiabatic calorimeter in temperature range from 80 to 390 K. The dependence of the molar heat capacity on temperature is given as a function of the reduced temperature (X) by polynomial equations, C P,m (J K−1 mol−1) = 204.75 + 81.421X − 23.828 X 2 + 12.044X 3 + 2.5442X 4 [X = (T − 132.5)/52.5] for the solid phase (80–185 K), C P,m (J K−1 mol−1) = 368.99 + 2.4199X + 1.0027X 2 + 0.43395X 3 [X = (T − 230)/35] for the glass state (195 − 265 K), and C P,m (J K−1 mol−1) = 415.01 + 21.992X − 0.24656X 2 + 0.57770X 3 [X = (T − 337.5)/52.5] for the liquid phase (285–390 K), respectively. According to the polynomial equations and thermodynamic relationship, the values of thermodynamic function of the BMIPF6 relative to 298.15 K were calculated in temperature range from 80 to 390 K with an interval of 5 K. The glass transition of BMIPF6 was measured to be 190.41 K, the enthalpy and entropy of the glass transition were determined to be ΔH g = 2.853 kJ mol−1 and ΔS g = 14.98 J K−1 mol−1, respectively. The results showed that the milting point of the BMIPF6 is 281.83 K, the enthalpy and entropy of phase transition were calculated to be ΔH m = 20.67 kJ mol−1 and ΔS m = 73.34 J K−1 mol−1.\n\nRestricted access\n\n## Thermodynamic properties of rubidium niobium tungsten oxide\n\nJournal of Thermal Analysis and Calorimetry\nAuthors:\nAleksandr Knyazev\n,\nMirosław Mączka\n,\nNataliya Kuznetsova\n,\nJerzy Hanuza\n, and\nAleksey Markin\n\n## Abstract\n\nIn the present work temperature dependence of heat capacity of rubidium niobium tungsten oxide has been measured first in the range from 7 to 395 K and then between 390 and 650 K, respectively, by precision adiabatic vacuum and dynamic calorimetry. The experimental data were used to calculate standard thermodynamic functions, namely the heat capacity\n\\documentclass{aastex} \\usepackage{amsbsy} \\usepackage{amsfonts} \\usepackage{amssymb} \\usepackage{bm} \\usepackage{mathrsfs} \\usepackage{pifont} \\usepackage{stmaryrd} \\usepackage{textcomp} \\usepackage{upgreek} \\usepackage{portland,xspace} \\usepackage{amsmath,amsxtra} \\pagestyle{empty} \\DeclareMathSizes{10}{9}{7}{6} \\begin{document} $$C_{\\text{p}}^{\\text{o}} (T),$$ \\end{document}\nenthalpy\n\\documentclass{aastex} \\usepackage{amsbsy} \\usepackage{amsfonts} \\usepackage{amssymb} \\usepackage{bm} \\usepackage{mathrsfs} \\usepackage{pifont} \\usepackage{stmaryrd} \\usepackage{textcomp} \\usepackage{upgreek} \\usepackage{portland,xspace} \\usepackage{amsmath,amsxtra} \\pagestyle{empty} \\DeclareMathSizes{10}{9}{7}{6} \\begin{document} $$H^{\\text{o}} ({\\rm T}) - H^{\\text{o}} (0)$$ \\end{document}\n, entropy\n\\documentclass{aastex} \\usepackage{amsbsy} \\usepackage{amsfonts} \\usepackage{amssymb} \\usepackage{bm} \\usepackage{mathrsfs} \\usepackage{pifont} \\usepackage{stmaryrd} \\usepackage{textcomp} \\usepackage{upgreek} \\usepackage{portland,xspace} \\usepackage{amsmath,amsxtra} \\pagestyle{empty} \\DeclareMathSizes{10}{9}{7}{6} \\begin{document} $$S^{\\text{o}} (T) - S^{\\text{o}} \\left( 0 \\right)$$ \\end{document}\n, and Gibbs function\n\\documentclass{aastex} \\usepackage{amsbsy} \\usepackage{amsfonts} \\usepackage{amssymb} \\usepackage{bm} \\usepackage{mathrsfs} \\usepackage{pifont} \\usepackage{stmaryrd} \\usepackage{textcomp} \\usepackage{upgreek} \\usepackage{portland,xspace} \\usepackage{amsmath,amsxtra} \\pagestyle{empty} \\DeclareMathSizes{10}{9}{7}{6} \\begin{document} $$G^{{^{\\text{o}} }} ({\\rm T}) - H^{{^{\\text{o}} }} (0)$$ \\end{document}\n, for the range from T→0 to 650 K. The high-temperature X-ray diffraction and the differential scanning calorimetry were used for the determination of temperature and decomposition products of RbNbWO6.\nRestricted access" ]
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https://ininet.org/language-specification-version-0-notice.html?page=34
[ "Language Specification Version 0\u0004\u0005 Notice\n\n Page 34/85 Date 29.01.2017 Size 3.2 Mb. #10878\n\n7.8Shift operators\n\nThe << and >> operators are used to perform bit shifting operations.\n\nshift-expression:\n\nFor an operation of the form x << count or x >> count, binary operator overload resolution (§7.2.4) is applied to select a specific operator implementation. The operands are converted to the parameter types of the selected operator, and the type of the result is the return type of the operator.\n\nWhen declaring an overloaded shift operator, the type of the first operand must always be the class or struct containing the operator declaration, and the type of the second operand must always be int.\n\nThe predefined shift operators are listed below.\n\n• Shift left:\n\nint operator <<(int x, int count);\nuint operator <<(uint x, int count);\nlong operator <<(long x, int count);\nulong operator <<(ulong x, int count);\n\nThe << operator shifts x left by a number of bits computed as described below.\n\nThe high-order bits outside the range of the result type of x are discarded, the remaining bits are shifted left, and the low-order empty bit positions are set to zero.\n\n• Shift right:\n\nint operator >>(int x, int count);\nuint operator >>(uint x, int count);\nlong operator >>(long x, int count);\nulong operator >>(ulong x, int count);\n\nThe >> operator shifts x right by a number of bits computed as described below.\n\nWhen x is of type int or long, the low-order bits of x are discarded, the remaining bits are shifted right, and the high-order empty bit positions are set to zero if x is non-negative and set to one if x is negative.\n\nWhen x is of type uint or ulong, the low-order bits of x are discarded, the remaining bits are shifted right, and the high-order empty bit positions are set to zero.\n\nFor the predefined operators, the number of bits to shift is computed as follows:\n\n• When the type of x is int or uint, the shift count is given by the low-order five bits of count. In other words, the shift count is computed from count & 0x1F.\n\n• When the type of x is long or ulong, the shift count is given by the low-order six bits of count. In other words, the shift count is computed from count & 0x3F.\n\nIf the resulting shift count is zero, the shift operators simply return the value of x.\n\nShift operations never cause overflows and produce the same results in checked and unchecked contexts.\n\nWhen the left operand of the >> operator is of a signed integral type, the operator performs an arithmetic shift right wherein the value of the most significant bit (the sign bit) of the operand is propagated to the high-order empty bit positions. When the left operand of the >> operator is of an unsigned integral type, the operator performs a logical shift right wherein high-order empty bit positions are always set to zero. To perform the opposite operation of that inferred from the operand type, explicit casts can be used. For example, if x is a variable of type int, the operation unchecked((int)((uint)x >> y)) performs a logical shift right of x.\n\n7.9Relational and type-testing operators\n\nThe ==, !=, <, >, <=, >=, is and as operators are called the relational and type-testing operators.\n\nrelational-expression:\nshift-expression\nrelational-expression < shift-expression\nrelational-expression > shift-expression\nrelational-expression <= shift-expression\nrelational-expression >= shift-expression\nrelational-expression is type\nrelational-expression as type\n\nequality-expression:\nrelational-expression\nequality-expression == relational-expression\nequality-expression != relational-expression\n\nThe is operator is described in §7.9.10 and the as operator is described in §7.9.11.\n\nThe ==, !=, <, >, <= and >= operators are comparison operators. For an operation of the form x op y, where op is a comparison operator, overload resolution (§7.2.4) is applied to select a specific operator implementation. The operands are converted to the parameter types of the selected operator, and the type of the result is the return type of the operator.\n\nThe predefined comparison operators are described in the following sections. All predefined comparison operators return a result of type bool, as described in the following table.\n\n Operation Result x == y true if x is equal to y, false otherwise x != y true if x is not equal to y, false otherwise x < y true if x is less than y, false otherwise x > y true if x is greater than y, false otherwise x <= y true if x is less than or equal to y, false otherwise x >= y true if x is greater than or equal to y, false otherwise\n\n7.9.1Integer comparison operators\n\nThe predefined integer comparison operators are:\n\nbool operator ==(int x, int y);\n\nbool operator ==(uint x, uint y);\nbool operator ==(long x, long y);\nbool operator ==(ulong x, ulong y);\n\nbool operator !=(int x, int y);\n\nbool operator !=(uint x, uint y);\nbool operator !=(long x, long y);\nbool operator !=(ulong x, ulong y);\n\nbool operator <(int x, int y);\n\nbool operator <(uint x, uint y);\nbool operator <(long x, long y);\nbool operator <(ulong x, ulong y);\n\nbool operator >(int x, int y);\n\nbool operator >(uint x, uint y);\nbool operator >(long x, long y);\nbool operator >(ulong x, ulong y);\n\nbool operator <=(int x, int y);\n\nbool operator <=(uint x, uint y);\nbool operator <=(long x, long y);\nbool operator <=(ulong x, ulong y);\n\nbool operator >=(int x, int y);\n\nbool operator >=(uint x, uint y);\nbool operator >=(long x, long y);\nbool operator >=(ulong x, ulong y);\n\nEach of these operators compares the numeric values of the two integer operands and returns a bool value that indicates whether the particular relation is true or false.\n\n7.9.2Floating-point comparison operators\n\nThe predefined floating-point comparison operators are:\n\nbool operator ==(float x, float y);\n\nbool operator ==(double x, double y);\n\nbool operator !=(float x, float y);\n\nbool operator !=(double x, double y);\n\nbool operator <(float x, float y);\n\nbool operator <(double x, double y);\n\nbool operator >(float x, float y);\n\nbool operator >(double x, double y);\n\nbool operator <=(float x, float y);\n\nbool operator <=(double x, double y);\n\nbool operator >=(float x, float y);\n\nbool operator >=(double x, double y);\n\nThe operators compare the operands according to the rules of the IEEE 754 standard:\n\n• If either operand is NaN, the result is false for all operators except !=, for which the result is true. For any two operands, x != y always produces the same result as !(x == y). However, when one or both operands are NaN, the <, >, <=, and >= operators do not produce the same results as the logical negation of the opposite operator. For example, if either of x and y is NaN, then x < y is false, but !(x >= y) is true.\n\n• When neither operand is NaN, the operators compare the values of the two floating-point operands with respect to the ordering\n\n–∞ < –max < ... < –min < –0.0 == +0.0 < +min < ... < +max < +∞\n\nwhere min and max are the smallest and largest positive finite values that can be represented in the given floating-point format. Notable effects of this ordering are:\n\n• Negative and positive zeros are considered equal.\n\n• A negative infinity is considered less than all other values, but equal to another negative infinity.\n\n• A positive infinity is considered greater than all other values, but equal to another positive infinity.\n\n7.9.3Decimal comparison operators\n\nThe predefined decimal comparison operators are:\n\nbool operator ==(decimal x, decimal y);\n\nbool operator !=(decimal x, decimal y);\n\nbool operator <(decimal x, decimal y);\n\nbool operator >(decimal x, decimal y);\n\nbool operator <=(decimal x, decimal y);\n\nbool operator >=(decimal x, decimal y);\n\nEach of these operators compares the numeric values of the two decimal operands and returns a bool value that indicates whether the particular relation is true or false. Each decimal comparison is equivalent to using the corresponding relational or equality operator of type System.Decimal.\n\n7.9.4Boolean equality operators\n\nThe predefined boolean equality operators are:\n\nbool operator ==(bool x, bool y);\n\nbool operator !=(bool x, bool y);\n\nThe result of == is true if both x and y are true or if both x and y are false. Otherwise, the result is false.\n\nThe result of != is false if both x and y are true or if both x and y are false. Otherwise, the result is true. When the operands are of type bool, the != operator produces the same result as the ^ operator.\n\n7.9.5Enumeration comparison operators\n\nEvery enumeration type implicitly provides the following predefined comparison operators:\n\nbool operator ==(E x, E y);\n\nbool operator !=(E x, E y);\n\nbool operator <(E x, E y);\n\nbool operator >(E x, E y);\n\nbool operator <=(E x, E y);\n\nbool operator >=(E x, E y);\n\nThe result of evaluating x op y, where x and y are expressions of an enumeration type E with an underlying type U, and op is one of the comparison operators, is exactly the same as evaluating ((U)x) op ((U)y). In other words, the enumeration type comparison operators simply compare the underlying integral values of the two operands.\n\n7.9.6Reference type equality operators\n\nThe predefined reference type equality operators are:\n\nbool operator ==(object x, object y);\n\nbool operator !=(object x, object y);\n\nThe operators return the result of comparing the two references for equality or non-equality.\n\nSince the predefined reference type equality operators accept operands of type object, they apply to all types that do not declare applicable operator == and operator != members. Conversely, any applicable user-defined equality operators effectively hide the predefined reference type equality operators.\n\nThe predefined reference type equality operators require one of the following:\n\n• Both operands are reference-type values or the value null. Furthermore, a standard implicit conversion (§6.3.1) exists from the type of either operand to the type of the other operand.\n\n• One operand is a value of type T where T is a type-parameter and the other operand is the value null. Furthermore T does not have the value type constraint.\n\nUnless one of these conditions are true, a compile-time error occurs. Notable implications of these rules are:\n\n• It is a compile-time error to use the predefined reference type equality operators to compare two references that are known to be different at compile-time. For example, if the compile-time types of the operands are two class types A and B, and if neither A nor B derives from the other, then it would be impossible for the two operands to reference the same object. Thus, the operation is considered a compile-time error.\n\n• The predefined reference type equality operators do not permit value type operands to be compared. Therefore, unless a struct type declares its own equality operators, it is not possible to compare values of that struct type.\n\n• The predefined reference type equality operators never cause boxing operations to occur for their operands. It would be meaningless to perform such boxing operations, since references to the newly allocated boxed instances would necessarily differ from all other references.\n\n• If an operand of a type parameter type T is compared to null, and the runtime type of T is a value type, the result of the comparison is false.\n\nThe following example checks whether an argument of an unconstrained type parameter type is null.\n\nclass C\n\n{\nvoid F(T x) {\nif (x == null) throw new ArgumentNullException();\n...\n}\n}\n\nThe x == null construct is permitted even though T could represent a value type, and the result is simply defined to be false when T is a value type.\n\nFor an operation of the form x == y or x != y, if any applicable operator == or operator != exists, the operator overload resolution (§7.2.4) rules will select that operator instead of the predefined reference type equality operator. However, it is always possible to select the predefined reference type equality operator by explicitly casting one or both of the operands to type object. The example\n\nusing System;\n\nclass Test\n{\nstatic void Main() {\nstring s = \"Test\";\nstring t = string.Copy(s);\nConsole.WriteLine(s == t);\nConsole.WriteLine((object)s == t);\nConsole.WriteLine(s == (object)t);\nConsole.WriteLine((object)s == (object)t);\n}\n}\n\nproduces the output\n\nTrue\nFalse\nFalse\nFalse\n\nThe s and t variables refer to two distinct string instances containing the same characters. The first comparison outputs True because the predefined string equality operator (§7.9.7) is selected when both operands are of type string. The remaining comparisons all output False because the predefined reference type equality operator is selected when one or both of the operands are of type object.\n\nNote that the above technique is not meaningful for value types. The example\n\nclass Test\n\n{\nstatic void Main() {\nint i = 123;\nint j = 123;\nSystem.Console.WriteLine((object)i == (object)j);\n}\n}\n\noutputs False because the casts create references to two separate instances of boxed int values.\n\n7.9.7String equality operators\n\nThe predefined string equality operators are:\n\nbool operator ==(string x, string y);\n\nbool operator !=(string x, string y);\n\nTwo string values are considered equal when one of the following is true:\n\n• Both values are null.\n\n• Both values are non-null references to string instances that have identical lengths and identical characters in each character position.\n\nThe string equality operators compare string values rather than string references. When two separate string instances contain the exact same sequence of characters, the values of the strings are equal, but the references are different. As described in §7.9.6, the reference type equality operators can be used to compare string references instead of string values.\n\n7.9.8Delegate equality operators\n\nEvery delegate type implicitly provides the following predefined comparison operators:\n\nbool operator ==(System.Delegate x, System.Delegate y);\n\nbool operator !=(System.Delegate x, System.Delegate y);\n\nTwo delegate instances are considered equal as follows:\n\n• If either of the delegate instances is null, they are equal if and only if both are null.\n\n• If the delegates have different runtime type they are never equal.\n\n• If both of the delegate instances have an invocation list (§15.1), those instances are equal if and only if their invocation lists are the same length, and each entry in one’s invocation list is equal (as defined below) to the corresponding entry, in order, in the other’s invocation list.\n\nThe following rules govern the equality of invocation list entries:\n\n• If two invocation list entries both refer to the same static method then the entries are equal.\n\n• If two invocation list entries both refer to the same non-static method on the same target object (as defined by the reference equality operators) then the entries are equal.\n\n• Invocation list entries produced from evaluation of semantically identical anonymous-function-expressions with the same (possibly empty) set of captured outer variable instances are permitted (but not required) to be equal.\n\n7.9.9Equality operators and null\n\nThe == and != operators permit one operand to be a value of a nullable type and the other to be the null literal, even if no predefined or user-defined operator (in unlifted or lifted form) exists for the operation.\n\nFor an operation of one of the forms\n\nx == null null == x x != null null != x\n\nwhere x is an expression of a nullable type, if operator overload resolution (§7.2.4) fails to find an applicable operator, the result is instead computed from the HasValue property of x. Specifically, the first two forms are translated into !x.HasValue, and last two forms are translated into x.HasValue.\n\n7.9.10The is operator\n\nThe is operator is used to dynamically check if the run-time type of an object is compatible with a given type. The result of the operation E is T, where E is an expression and T is a type, is a boolean value indicating whether E can successfully be converted to type T by a reference conversion, a boxing conversion, or an unboxing conversion. The operation is evaluated as follows, after type arguments have been substituted for all type parameters:\n\n• If E is an anonymous function, a compile time error occurs\n\n• If E is a method group or the null literal, of if the type of E is a reference type or a nullable type and the value of E is null, the result is false.\n\n• Otherwise, let D represent the dynamic type of E as follows:\n\n• If the type of E is a reference type, D is the run-time type of the instance reference by E.\n\n• If the type of E is a nullable type, D is the underlying type of that nullable type.\n\n• If the type of E is a non-nullable value type, D is the type of E.\n\n• The result of the operation depends on D and T as follows:\n\n• If T is a reference type, the result is true if D and T are the same type, if D is a reference type and an implicit reference conversion from D to T exists, or if D is a value type and a boxing conversion from D to T exists.\n\n• If T is a nullable type, the result is true if D is the underlying type of T.\n\n• If T is a non-nullable value type, the result is true if D and T are the same type.\n\n• Otherwise, the result is false.\n\nNote that user defined conversions, are not considered by the is operator.\n\n7.9.11The as operator\n\nThe as operator is used to explicitly convert a value to a given reference type or nullable type. Unlike a cast expression (§7.6.6), the as operator never throws an exception. Instead, if the indicated conversion is not possible, the resulting value is null.\n\nIn an operation of the form E as T, E must be an expression and T must be a reference type, a type parameter known to be a reference type, or a nullable type. Furthermore, at least one of the following must be true, or otherwise a compile-time error occurs:\n\n• An identity (§6.1.1), implicit reference (§6.1.6), boxing (§6.1.7), explicit reference (§6.2.4), or unboxing (§6.2.5) conversion exists from the type of E to T.\n\n• The type of E or T is an open type.\n\n• E is the null literal.\n\nThe operation E as T produces the same result as\n\nE is T ? (T)(E) : (T)null\n\n• except that E is only evaluated once. The compiler can be expected to optimize E as T to perform at most one dynamic type check as opposed to the two dynamic type checks implied by the expansion above.\n\nNote that some conversions, such as user defined conversions, are not possible with the as operator and should instead be performed using cast expressions.\n\nIn the example\n\nclass X\n{\n\npublic string F(object o) {\n\nreturn o as string; // OK, string is a reference type\n}\n\npublic T G(object o) where T: Attribute {\n\nreturn o as T; // Ok, T has a class constraint\n}\n\npublic U H(object o) {\n\nreturn o as U; // Error, U is unconstrained\n}\n}\n\nthe type parameter T of G is known to be a reference type, because it has the class constraint. The type parameter U of H is not however; hence the use of the as operator in H is disallowed." ]
[ null ]
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https://mathhelpboards.com/threads/automorphism-of-order-2-fixing-just-identity-prove-that-g-is-abelian.2833/
[ "# automorphism of order 2 fixing just identity. Prove that G is abelian.\n\n#### caffeinemachine\n\n##### Well-known member\nMHB Math Scholar\nLet $G$ be a finite group, $T$ an automorphism of $G$ with the property that $T(x)=x$ if and only if $x=e$. Suppose further that $T^2=I$, that is, $T(T(x))=x$ for all $x\\in G$. Show that $G$ is abelian.\n\nI approached this problem using the permutation representation afforded by $T$ on $G$. Its easy to deduce that the cycle representation of the permutation of $G$ caused by $T$ has $(n-1)/2$ disjoint transpositions, where $n=|G|$. We know, from this, that $n$ is odd but so what? I am not able to exploit the homomorphism property of $T$.\n\n#### Deveno\n\n##### Well-known member\nMHB Math Scholar\nhint: (i don't want to spoil your fun because this is a beautiful theorem)\n\ndefine $f:G \\to G$ by:\n\n$f(g) = g^{-1}T(g)$\n\nshow $f$ is injective (and thus bijective).\n\nnow...what is $T\\circ f(g)$?" ]
[ null ]
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https://www.farmrio.com/collections/midi-dresses-sale
[ "# Midi Dresses Sale", null, "Color\n• MULTI\nSize\n• XS\n• S\n• M\n3 left\n• L\n2 left\n• XL\n2 left\nColor\n• WHITE\nSize\n• XS\n• S\n• M\n• L\n1 left\n• XL\nColor\n• MULTI\nSize\n• XS\n• S\n• M\n• L\n• XL\nColor\n• MULTI\nSize\n• XS\n• S\n• M\n2 left\n• L\n• XL\n2 left\nColor\n• WHITE\nSize\n• XS\n• S\n• M\n• L\n• XL\nColor\n• MULTI\nSize\n• XS\n3 left\n• S\n• M\n2 left\n• L\n3 left\n• XL\nColor\n• Yellow MULTI\nSize\n• XS\n• S\n3 left\n• M\n1 left\n• L\n3 left\n• XL\n3 left\nColor\n• Red MULTI\nSize\n• XS\n• S\n• M\n2 left\n• L\n• XL\nColor\n• MULTI\nSize\n• XS\n• S\n• M\n• L\n• XL\nColor\n• OFF-WHITE\nSize\n• XS\n• S\n• M\n• L\n• XL\nColor\n• PINK MULTI\nSize\n• XS\n• S\n• M\n• L\n• XL\nColor\n• Red MULTI\nSize\n• XS\n• S\n• M\n• L\n1 left\n• XL\nColor\n• MULTI\nSize\n• XS\n• S\n• M\n• L\n• XL\nColor\n• MULTI\nSize\n• XS\n• S\n• M\n• L\n• XL\nColor\n• BLUE MULTI\nSize\n• XS\n• S\n• M\n• L\n• XL\nColor\n• MULTI\nSize\n• XS\n• S\n• M\n• L\n• XL\nColor\n• MULTI\nSize\n• XS\n• S\n• M\n1 left\n• L\n• XL\nColor\n• Red MULTI\nSize\n• XS\n• S\n• M\n1 left\n• L\n• XL" ]
[ null, "https://cdn.shopify.com/s/files/1/0077/6673/6963/collections/banner_subcategory_sale_final_1920x.jpg", null ]
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https://nl.mathworks.com/help/images/ref/deconvblind.html
[ "# deconvblind\n\nDeblur image using blind deconvolution\n\n## Syntax\n\n``[J,psfr] = deconvblind(I,psfi)``\n``[J,psfr] = deconvblind(I,psfi,iter)``\n``[J,psfr] = deconvblind(I,psfi,iter,dampar)``\n``[J,psfr] = deconvblind(I,psfi,iter,dampar,weight)``\n``[J,psfr] = deconvblind(I,psfi,iter,dampar,weight,readout)``\n``[J,psfr] = deconvblind(___,fun)``\n\n## Description\n\n````[J,psfr] = deconvblind(I,psfi)` deconvolves image `I` using the maximum likelihood algorithm and an initial estimate of the point-spread function (PSF), `psfi`. The `deconvblind` function returns both the deblurred image `J` and a restored PSF, `psfr`.To improve the restoration, `deconvblind` supports several optional parameters, described below. Use `[]` as a placeholder if you do not specify an intermediate parameter.```\n````[J,psfr] = deconvblind(I,psfi,iter)` specifies the number of iterations, `iter`.```\n````[J,psfr] = deconvblind(I,psfi,iter,dampar)` controls noise amplification by suppressing iterations for pixels that deviate a small amount compared to the noise, specified by the damping threshold `dampar`. By default, no damping occurs.```\n\nexample\n\n````[J,psfr] = deconvblind(I,psfi,iter,dampar,weight)` specifies which pixels in the input image `I` are considered in the restoration. The value of an element in the `weight` array determines how much the pixel at the corresponding position in the input image is considered. For example, to exclude a pixel from consideration, assign it a value of `0` in the `weight` array. You can adjust the weight value assigned to each pixel according to the amount of flat-field correction.```\n````[J,psfr] = deconvblind(I,psfi,iter,dampar,weight,readout)` specifies the additive noise (such as background and foreground noise) and the variance of the read-out camera noise, `readout`.```\n````[J,psfr] = deconvblind(___,fun)`, where `fun` is a handle to a function that describes additional constraints on the PSF. `fun` is called at the end of each iteration. For more information about function handles, see Create Function Handle.```\n\n## Examples\n\ncollapse all\n\nCreate a sample image with noise.\n\n```% Set the random number generator back to its default settings for % consistency in results. rng default; I = checkerboard(8); PSF = fspecial('gaussian',7,10); V = .0001; BlurredNoisy = imnoise(imfilter(I,PSF),'gaussian',0,V);```\n\nCreate a weight array to specify which pixels are included in processing.\n\n```WT = zeros(size(I)); WT(5:end-4,5:end-4) = 1; INITPSF = ones(size(PSF));```\n\nPerform blind deconvolution.\n\n`[J P] = deconvblind(BlurredNoisy,INITPSF,20,10*sqrt(V),WT);`\n\nDisplay the results.\n\n```subplot(221);imshow(BlurredNoisy); title('A = Blurred and Noisy'); subplot(222);imshow(PSF,[]); title('True PSF'); subplot(223);imshow(J); title('Deblurred Image'); subplot(224);imshow(P,[]); title('Recovered PSF');```", null, "## Input Arguments\n\ncollapse all\n\nBlurry image, specified as a numeric array of any dimension. You can also specify the image as a cell array to enable interrupted iterations. For more information, see Tips.\n\nData Types: `single` | `double` | `int16` | `uint8` | `uint16`\n\nInitial estimate of PSF, specified as a numeric array. The PSF restoration is affected strongly by the size of the initial guess `psfi` and less by the values it contains. For this reason, specify an array of `1`s as your `psfi`.\n\nYou can also specify `psfi` as a cell array to enable interrupted iterations. For more information, see Tips.\n\nData Types: `single` | `double` | `int16` | `uint8` | `uint16`\n\nNumber of iterations, specified as a positive integer.\n\nData Types: `double`\n\nThreshold for damping, specified as a numeric scalar. Damping occurs for pixels whose deviation between iterations is less than the threshold. `dampar` has the same data type as `I`.\n\nWeight value of each pixel, specified as a numeric array with values in the range [0, 1]. `weight` has the same size as the input image, `I`. By default, all elements in `weight` have the value `1`, so all pixels are considered equally in the restoration.\n\nData Types: `double`\n\nNoise, specified as a numeric scalar or numeric array. The value of `readout` corresponds to the additive noise (such as noise from the foreground and background) and the variance of the read-out camera noise. `readout` has the same data type as `I`.\n\nFunction handle, specified as a handle. `fun` must accept the PSF as its first argument. The function must return one argument: a PSF that is the same size as the original PSF and that satisfies the positivity and normalization constraints.\n\n## Output Arguments\n\ncollapse all\n\nDeblurred image, returned as a numeric array or a 1-by-4 cell array. `J` (or `J{1}` when `J` is a cell array) has the same data type as `I`. For more information about returning `J` as a cell array for interrupted iterations, see Tips.\n\nRestored PSF, returned as an array of positive numbers or a 1-by-4 cell array. `psfr` has the same size as the initial estimate of the PSF, `psfi`, and it is normalized so the sum of elements is 1. For more information about returning `psfr` as a cell array for interrupted iterations, see Tips.\n\nData Types: `double`\n\n## Tips\n\n• You can use `deconvblind` to perform a deconvolution that starts where a previous deconvolution stopped. To use this feature, pass the input image `I` and the initial guess at the PSF, `psfi`, as cell arrays: `{I}` and `{psfi}`. When you do, the `deconvblind` function returns the output image `J` and the restored point-spread function, `psfr`, as cell arrays, which can then be passed as the input arrays into the next `deconvblind` call. The output cell array `J` contains four elements:\n\n`J{1}` contains `I`, the original image.\n\n`J{2}` contains the result of the last iteration.\n\n`J{3}` contains the result of the next-to-last iteration.\n\n`J{4}` is an array generated by the iterative algorithm.\n\n• The output image `J` could exhibit ringing introduced by the discrete Fourier transform used in the algorithm. To reduce the ringing, use `I = edgetaper(I,psfi)` before calling `deconvblind`.\n\n D.S.C. Biggs and M. Andrews, Acceleration of iterative image restoration algorithms, Applied Optics, Vol. 36, No. 8, 1997.\n\n R.J. Hanisch, R.L. White, and R.L. Gilliland, Deconvolutions of Hubble Space Telescope Images and Spectra, Deconvolution of Images and Spectra, Ed. P.A. Jansson, 2nd ed., Academic Press, CA, 1997.\n\n Timothy J. Holmes, et al, Light Microscopic Images Reconstructed by Maximum Likelihood Deconvolution, Handbook of Biological Confocal Microscopy, Ed. James B. Pawley, Plenum Press, New York, 1995." ]
[ null, "https://nl.mathworks.com/help/examples/images/win64/DeblurAnImageUsingBlindDeconvolutionExample_01.png", null ]
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http://st.tongbuedu.com/counts-33114.html
[ "", null, "同步学习网是一个提供了大量的中考真题下载,高考真题下载,中考试题解析,高考试题解析,各个学习段的模拟试题下载的网站,我们希望通过在线做题的方式能够切实的帮助广大学生查缺补漏,提高学习成绩!", null, "1. 已知全集", null, ",集合", null, "", null, ",则", null, "A.{0}                          B.{-3-4}                  C.{-4-2}                  D.", null, "2. 已知集合", null, "", null, ",则", null, "A.", null, "B.", null, "C.", null, "D.", null, "3. 设函数", null, "在区间(12)内有零点,则实数a的取值范围是(\n\nA.", null, "", null, "B.", null, "", null, "C.", null, "", null, "D.", null, "", null, "5. 下面几种推理中是演绎推理的序号为(\n\nA. 由金、银、铜、铁可导电,猜想:金属都可导电;\n\n6. 在下列命题中,真命题是(\n\nA.x=2,x23x+2=0的否命题;         B.b=3,b2=9的逆命题;\n\nC. ac>bc,a>b;                          D.相似三角形的对应角相等的逆否命题\n\n7. 函数", null, "的图象大致是", null, "8. 已知命题p", null, "x1x2", null, "R(f(x2)", null, "f(x1)(x2", null, "x1)0,则", null, "p是(\n\nA.", null, "x1x2", null, "R(f(x2)", null, "f(x1)(x2", null, "x1)0\n\nB.", null, "x1x2", null, "R(f(x2)", null, "f(x1)(x2", null, "x1)0\n\nC.", null, "x1x2", null, "R(f(x2)", null, "f(x1)(x2", null, "x1)<0\n\nD.", null, "x1x2", null, "R(f(x2)", null, "f(x1)(x2", null, "x1)<0\n\nA.", null, "B.", null, "C.", null, "D.", null, "10. 定义在", null, "上的奇函数", null, ",当", null, "时,", null, ",则关于", null, "的函数", null, "的所有零点之和为(\n\nA.", null, "B.", null, "C.", null, "D.", null, "11. 若函数", null, "的单调递增区间是", null, ",则", null, "_________.\n\n12. 曲线", null, "在点", null, "处的切线斜率为        .\n\n13. 已知函数", null, ",则f (4) =       .\n\n14. 把数列", null, "依次按第一个括号一个数,第二个括号两个数,第三个括号三个数,第四个括号四个数,第五个括号一个数,第六个括号两个数,……,循环下去,如:(3),(57),(91113),(15171921),(23),(2527),……,则第104个括号内各数字之和为________.\n\n15.", null, "是定义在", null, "上且周期为2的函数,在区间", null, "上,", null, "其中", null, ".", null, ",则", null, "的值为___________.\n\n16.(本题满分8分)已知命题", null, ",命题", null, "(其中m > 0),且", null, "的必要条件,求实数m的取值范围。\n\n17.(本题满分8分)\n\nx=1x=2是函数f(x)=alnx+bx2+x的两个极值点\n\n(1)a,b的值;\n\n(2)f(x)的单调区间。\n\n18.(本题满分10分)\n\n19.(本题满分12分)\n\n20.(本题满分13分)", null, "" ]
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{"ft_lang_label":"__label__zh","ft_lang_prob":0.9551222,"math_prob":0.99930847,"size":1392,"snap":"2020-24-2020-29","text_gpt3_token_len":1258,"char_repetition_ratio":0.12896253,"word_repetition_ratio":0.03508772,"special_character_ratio":0.44899425,"punctuation_ratio":0.20858896,"nsfw_num_words":1,"has_unicode_error":false,"math_prob_llama3":0.957966,"pos_list":[0,1,2,3,4,5,6,7,8,9,10,11,12,13,14,15,16,17,18,19,20,21,22,23,24,25,26,27,28,29,30,31,32,33,34,35,36,37,38,39,40,41,42,43,44,45,46,47,48,49,50,51,52,53,54,55,56,57,58,59,60,61,62,63,64,65,66,67,68,69,70,71,72,73,74,75,76,77,78,79,80,81,82,83,84,85,86,87,88,89,90,91,92,93,94,95,96,97,98,99,100,101,102,103,104,105,106,107,108,109,110,111,112,113,114,115,116,117,118,119,120,121,122,123,124,125,126,127,128,129,130,131,132,133,134,135,136,137,138,139,140,141,142,143,144,145,146,147,148,149,150,151,152,153,154,155,156],"im_url_duplicate_count":[null,null,null,null,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,6,null,10,null,null,null,null,null,2,null,4,null,10,null,null,null,null,null,6,null,10,null,null,null,null,null,4,null,10,null,null,null,null,null,6,null,10,null,null,null,null,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,2,null,null,null],"WARC_HEADER":"{\"WARC-Type\":\"response\",\"WARC-Date\":\"2020-07-07T08:49:03Z\",\"WARC-Record-ID\":\"<urn:uuid:5ae82269-937b-43ee-8772-2a45d1e633f0>\",\"Content-Length\":\"140009\",\"Content-Type\":\"application/http; msgtype=response\",\"WARC-Warcinfo-ID\":\"<urn:uuid:3dd1156b-008e-4565-8751-ba0e229b13bf>\",\"WARC-Concurrent-To\":\"<urn:uuid:46fbb4ed-84d2-4d97-bc3d-0f8b68e3bb56>\",\"WARC-IP-Address\":\"125.76.233.118\",\"WARC-Target-URI\":\"http://st.tongbuedu.com/counts-33114.html\",\"WARC-Payload-Digest\":\"sha1:KMEAZO4DESRHUF663MCKNY3TLAVHLWNY\",\"WARC-Block-Digest\":\"sha1:6KHNAIHXWXVOCGKFGBQ32S4IRYMAPXFT\",\"WARC-Identified-Payload-Type\":\"application/xhtml+xml\",\"warc_filename\":\"/cc_download/warc_2020/CC-MAIN-2020-29/CC-MAIN-2020-29_segments_1593655891884.11_warc_CC-MAIN-20200707080206-20200707110206-00157.warc.gz\"}"}
http://tool.114la.com/bus/shenzhen/s_R_65.html
[ "8路 德兴花园 --- 火车站西广场\n\n1.\n\n2.\n\n3.\n\n4.\n\n5.\n\n6.\n\n7.\n\n8.\n\n9.\n\n10.\n\n11.\n\n12.\n\n13.\n\n14.\n\n15.\n\n16.\n\n17.\n\n18.\n\n19.\n\n20.\n\n9路 德兴花园总站 --- 皇岗口岸\n\n1.\n\n2.\n\n3.\n\n4.\n\n5.\n\n6.\n\n7.\n\n8.\n\n9.\n\n10.\n\n11.\n\n12.\n\n13.\n\n14.\n\n15.\n\n16.\n\n17.\n\n18.\n\n19.\n\n20.\n\n21.\n\nN16路 德兴花园总站 --- 火车站\n\n1.\n\n2.\n\n3.\n\n4.\n\n5.\n\n6.\n\n7.\n\n8.\n\n9.\n\n10.\n\n11.\n\n12.\n\n13.\n\n14.\n\n15.\n\n16.\n\n17.\n\n18.\n\n19.\n\n20.\n\n57路 鹤围村 --- 长岭(总站)\n\n1.\n\n2.\n\n3.\n\n4.\n\n5.\n\n6.\n\n7.\n\n8.\n\n9.\n\n10.\n\n11.\n\n12.\n\n13.\n\n14.\n\n15.\n\n16.\n\n17.\n\n18.\n\n19.\n\n20.\n\n21.\n\n22.\n\n23.\n\n24.\n\n25.\n\n26.\n\n27.\n\n58路 南头火车西站 --- 德兴花园总站\n\n1.\n\n2.\n\n3.\n\n4.\n\n5.\n\n6.\n\n7.\n\n8.\n\n9.\n\n10.\n\n11.\n\n12.\n\n13.\n\n14.\n\n15.\n\n16.\n\n17.\n\n18.\n\n19.\n\n20.\n\n21.\n\n22.\n\n23.\n\n24.\n\n25.\n\n26.\n\n27.\n\n28.\n\n29.\n\n30.\n\n31.\n\n32.\n\n33.\n\n34.\n\n35.\n\n36.\n\n37.\n\n38.\n\n39.\n\n40.\n\n41.\n\n42.\n\n43.\n\n44.\n\n45.\n\n46.\n\n47.\n\n61路 德兴花园 --- 火车站西广场\n\n1.\n\n2.\n\n3.\n\n4.\n\n5.\n\n6.\n\n7.\n\n8.\n\n9.\n\n10.\n\n11.\n\n12.\n\n13.\n\n14.\n\n15.\n\n16.\n\n17.\n\n18.\n\n19.\n\n79路 东角头总站 --- 清水河总站\n\n1.\n\n2.\n\n3.\n\n4.\n\n5.\n\n6.\n\n7.\n\n8.\n\n9.\n\n10.\n\n11.\n\n12.\n\n13.\n\n14.\n\n15.\n\n16.\n\n17.\n\n18.\n\n19.\n\n20.\n\n21.\n\n22.\n\n23.\n\n24.\n\n25.\n\n26.\n\n27.\n\n28.\n\n29.\n\n30.\n\n31.\n\n32.\n\n33.\n\n34.\n\n35.\n\n36.\n\n37.\n\n38.\n\n39.\n\n40.\n\n41.\n\n42.\n\n43.\n\n44.\n\n45.\n\n46.\n\n47.\n\n48.\n\n85路 盐田检查站 --- 德兴花园\n\n1.\n\n2.\n\n3.\n\n4.\n\n5.\n\n6.\n\n7.\n\n8.\n\n9.\n\n10.\n\n11.\n\n12.\n\n13.\n\n14.\n\n15.\n\n16.\n\n17.\n\n18.\n\n19.\n\n20.\n\n21.\n\n22.\n\n23.\n\n24.\n\n25.\n\n26.\n\n27.\n\n28.\n\n29.\n\n30.\n\n31.\n\n32.\n\n33.\n\n34.\n\n35.\n\n36.\n\n37.\n\n38.\n\n39.\n\n40.\n\n41.\n\nK105路 蛇口码头 --- 德兴花园总站\n\n1.\n\n2.\n\n3.\n\n4.\n\n5.\n\n6.\n\n7.\n\n8.\n\n9.\n\n10.\n\n11.\n\n12.\n\n13.\n\n14.\n\n15.\n\n16.\n\n17.\n\n18.\n\n19.\n\n20.\n\n21.\n\n22.\n\n23.\n\n24.\n\n25.\n\n26.\n\n27.\n\n28.\n\n29.\n\n30.\n\n31.\n\n32.\n\n33.\n\n34.\n\n35.\n\n36.\n\n37.\n\n38.\n\n39.\n\n40.\n\n123路 深圳湾口岸 --- 德兴花园\n\n1.\n\n2.\n\n3.\n\n4.\n\n5.\n\n6.\n\n7.\n\n8.\n\n9.\n\n10.\n\n11.\n\n12.\n\n13.\n\n14.\n\n15.\n\n16.\n\n17.\n\n18.\n\n19.\n\n20.\n\n21.\n\n22.\n\n23.\n\n24.\n\n25.\n\n26.\n\n27.\n\n28.\n\n29.\n\n30.\n\n31.\n\n32.\n\n33.\n\n34.\n\n35.\n\n36.\n\n37.\n\n38.\n\n39.\n\n40.\n\n41.\n\n42.\n\n43.\n\n219路 石厦南总站 --- 清水河总站\n\n1.\n\n2.\n\n3.\n\n4.\n\n5.\n\n6.\n\n7.\n\n8.\n\n9.\n\n10.\n\n11.\n\n12.\n\n13.\n\n14.\n\n15.\n\n16.\n\n17.\n\n18.\n\n19.\n\n20.\n\n21.\n\n22.\n\n23.\n\n24.\n\n25.\n\n26.\n\n27.\n\n28.\n\n29.\n\n30.\n\n31.\n\n32.\n\n33.\n\n222路 德兴花园 --- 世界之窗总站\n\n1.\n\n2.\n\n3.\n\n4.\n\n5.\n\n6.\n\n7.\n\n8.\n\n9.\n\n10.\n\n11.\n\n12.\n\n13.\n\n14.\n\n15.\n\n16.\n\n17.\n\n18.\n\n19.\n\n20.\n\n21.\n\n22.\n\n23.\n\n24.\n\n25.\n\n26.\n\n27.\n\n28.\n\n29.\n\n30.\n\n31.\n\n32.\n\n33.\n\n34.\n\n35.\n\n36.\n\n37.\n\n385路 福田保税区 --- 元平学校\n\n1.\n\n2.\n\n3.\n\n4.\n\n5.\n\n6.\n\n7.\n\n8.\n\n9.\n\n10.\n\n11.\n\n12.\n\n13.\n\n14.\n\n15.\n\n16.\n\n17.\n\n18.\n\n19.\n\n20.\n\n21.\n\n22.\n\n23.\n\n24.\n\n25.\n\n26.\n\n27.\n\n28.\n\n29.\n\nB616路 清水河 --- 紫星花园\n\n1.\n\n2.\n\n3.\n\n4.\n\n5.\n\n6.\n\n### 热门车站\n\n 深圳 晴转多云 32℃ ~25℃" ]
[ null ]
{"ft_lang_label":"__label__zh","ft_lang_prob":0.9724755,"math_prob":0.92402774,"size":1098,"snap":"2020-34-2020-40","text_gpt3_token_len":1034,"char_repetition_ratio":0.22029251,"word_repetition_ratio":0.0,"special_character_ratio":0.5364299,"punctuation_ratio":0.17786561,"nsfw_num_words":0,"has_unicode_error":false,"math_prob_llama3":0.9653503,"pos_list":[0],"im_url_duplicate_count":[null],"WARC_HEADER":"{\"WARC-Type\":\"response\",\"WARC-Date\":\"2020-10-01T23:52:53Z\",\"WARC-Record-ID\":\"<urn:uuid:81769b66-907a-4e9e-b592-eea4fc7a27e9>\",\"Content-Length\":\"56694\",\"Content-Type\":\"application/http; msgtype=response\",\"WARC-Warcinfo-ID\":\"<urn:uuid:8f1be624-4256-4641-a872-62db49240017>\",\"WARC-Concurrent-To\":\"<urn:uuid:e0afcf8c-a1ef-477c-94d6-1bae097aaea4>\",\"WARC-IP-Address\":\"39.107.146.250\",\"WARC-Target-URI\":\"http://tool.114la.com/bus/shenzhen/s_R_65.html\",\"WARC-Payload-Digest\":\"sha1:OAMNAH4L7Z3LLPR7TJYX4OFLX3QKIXQO\",\"WARC-Block-Digest\":\"sha1:5GKOTPHNC72N5NY3XJY5NPRZRNQFET4X\",\"WARC-Identified-Payload-Type\":\"application/xhtml+xml\",\"warc_filename\":\"/cc_download/warc_2020/CC-MAIN-2020-40/CC-MAIN-2020-40_segments_1600402132335.99_warc_CC-MAIN-20201001210429-20201002000429-00014.warc.gz\"}"}
http://lib.mexmat.ru/books/127579
[ "Электронная библиотека Попечительского советамеханико-математического факультета Московского государственного университета\n Главная    Ex Libris    Книги    Журналы    Статьи    Серии    Каталог    Wanted    Загрузка    ХудЛит    Справка    Поиск по индексам    Поиск    Форум", null, "Авторизация", null, "Поиск по указателям", null, "", null, "", null, "", null, "", null, "Richard Haberman — Mathematical models", null, "Читать книгубесплатно\n\nСкачать книгу с нашего сайта нельзя\n\nОбсудите книгу на научном форуме", null, "Нашли опечатку?\nВыделите ее мышкой и нажмите Ctrl+Enter\n\nНазвание: Mathematical models\n\nАвтор: Richard Haberman\n\nАннотация:\n\nMathematics is a grand subject in the way it can be applied to various problems in science and engineering. To use mathematics, one needs to understand the physical context. The author uses mathematical techniques along with observations and experiments to give an in-depth look at models for mechanical vibrations, population dynamics, and traffic flow. Equal emphasis is placed on the mathematical formulation of the problem and the interpretation of the results. In the sections on mechanical vibrations and population dynamics, the author emphasizes the nonlinear aspects of ordinary differential equations and develops the concepts of equilibrium solutions and their stability. He introduces phase plane methods for the nonlinear pendulum and for predator-prey and competing species models.\n\nЯзык:", null, "Статус предметного указателя: Неизвестно\n\ned2k: ed2k stats\n\nГод издания: 1987\n\nКоличество страниц: 421\n\nДобавлена в каталог: 29.11.2013\n\nОперации: Положить на полку | Скопировать ссылку для форума | Скопировать ID", null, "Предметный указатель", null, "Реклама", null, "", null, "", null, "", null, "", null, "© Электронная библиотека попечительского совета мехмата МГУ, 2004-2019", null, "|", null, "|", null, "О проекте" ]
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http://fenixdocs.com/index.php/False
[ "This wiki is out of date, use the continuation of this wiki instead\n\n# FALSE\n\n(Redirected from False)\n\n## Definition\n\nINT FALSE\n\nFalse is a constant integer, equal to the value 0. It is used to state that something is false and not true.\n\nChecking whether a variable is false, is the same as checking if it's zero. In older versions of Fenix, it was the same as checking whether a variable is even. This has been changed because it's more commonly used.\n\n## Example\n\n```Program example;\nPrivate\nint b = false;\nBegin\n\n// comparison with the constant FALSE\nif(b == false)\nsay(\"b was FALSE! so b==0\");\nelse\nsay(\"b was not FALSE! so b!=0\");\nend\n\n// checking the integer itself\nif(!b)\nsay(\"b was false! so b==0\");\nelse\nsay(\"b was true! so b!=0\");\nend\n\nLoop\nframe;\nEnd\n\nEnd\n```" ]
[ null ]
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https://www.scirp.org/journal/paperinformation.aspx?paperid=64342
[ "Thermal Decomposition and a Kinetic Study of Poly(Para-Substituted Styrene)s\n\nAbstract\n\nThe thermal decompositions of polystyrene (PS), poly(p-methyl styrene) (PMS), poly(p-bromo styrene) (PBrS), and poly(p-chloro styrene) (PClS) were investigated through thermogravimetric analysis (TGA). For this aim, Flynn-Wall-Ozawa method was applied to derivative thermogravimetric (DTG) curves. Continuous distribution kinetics was employed with a stoichiometric kernel to determine the rate coefficients for decomposition reactions. TGA data for the polymers were investigated by nonlinear fitting procedures that yielded activation energies and frequency factors for the combined chemical reactions. The reaction order values of PS derivatives are just about 1 in the nonisothermal decomposition process. Ea values for PS, PMS, and PClS increase with % conversion individually as they decrease in the order of PS/PMS/PClS which is consistent with the molecular weight increase. On the other hand, PBrS has the highest activation energy. Also its activation energy decreases with the % conversion. Thus it is suggested that PBrS degrades with somehow different mechanism.\n\nKeywords\n\nShare and Cite:\n\nŞenocak, A. , Alkan, C. and Karadağ, A. (2016) Thermal Decomposition and a Kinetic Study of Poly(Para-Substituted Styrene)s. American Journal of Analytical Chemistry, 7, 246-253. doi: 10.4236/ajac.2016.73021.\n\nReceived 8 January 2016; accepted 6 March 2016; published 9 March 2016", null, "1. Introduction\n\nOverview and Background\n\nPolystyrene (PS) is a thermoplastic material having so many using areas like home appliances, construction, packaging, and other industrial applications. It is cheap and easily processable. However, it has restricted mechanical properties and is very flammable. Due to its advantages, different kinds of derivatives have been produced. To enhance PS’s competition capacity in engineering resin applications, to protect its properties like thermal stability is crucial. In this work, poly(para-substituted styrene)s have been investigated using nonisothermal degradation kinetics.\n\nUsually, polystyrene, exposed to thermal degradation, softens and melts at 160˚C. Up to 275˚C, volatility of fused polymer with high molecular weight is lowered and this molten version of the polymer keeps its stability. However, polymer fragments with lower molecular weight have higher volatility form and the decomposition is completed after the temperature reaches to 470˚C . In addition, molecular weight rapidly decreases in the beginning of the degradation and the rate of this decline slows down at later stages . This case applies to not only polystyrene but also some of its substituted derivatives. Literature suggests that styrene is the major decomposition product in the complete depolymerization taking place below 500˚C. After this temperature, further fragmentation carries out and the oligomers turn into gaseous species having lower molecular weight .\n\nFor kinetic interpretations in this study, poly(p-substituted styrene)s’ thermal degradation data were collected in nitrogen atmosphere. With intent to understand the origin of the equations to determine the kinetic parameters, a brief derivation is included below. For a simple degradation step, the rate of conversion is described as follows:", null, "(1)\n\nIn this equation, a is the degree of chemical conversion, n is the reaction order, k is the reaction rate, and d/dt is the time derivative. A relation between the reaction rate and temperature is usually assumed:", null, "(2)\n\nHere, Z is the preexponential constant, R is the ideal gas constant, Ea is the activation energy, and T is the temperature in Kelvin. This is known as the Arrhenius relationship. Equation (3) is attained by combining Equations (1) and (2):", null, "(3)\n\nIn order to obtain a precise solution, one of three variables in Equation (3) either may be held as a constant or may be related to another variable when the reaction process is carried out. In this way, the reaction process can be defined. For years, a differential or an integral version of Equation (3) has been applied by researchers to calculate the kinetic parameters consisting of Ea, Z, and n. While some of the techniques rely on the functional version of conversion, f(a), in the literature - , activation energy values have also been indicated to be estimated without any preceding knowledge of the functional form.\n\nThere is a linear relationship between temperature and time as following:", null, "(4)\n\nand", null, "(5)\n\nwhere T is the current temperature, T(0) is the initial temperature, t is time, and b is the scanning rate. Combining Equation (3) with Equation (5) yields:", null, "(6)\n\nThis equation shows the theoretical shape of the TGA curve that is determined in the Kinetics Software. Logarithmic form of the equation is", null, "(7)\n\nA multilinear regression is performed by using ln(bda/dT), −l/RT, and ln(1−a) which have been appraised from TGA curve, and solving for Z, Ea, and n.\n\nKinetic parameters of degradation are determined upon the observations of Ozawa and Flynn and Wall . Their statements involve an approximation that has been discussed by Doyle and by Flynn and Wall , among others. This approximation comprises the so-called “exponential integral” and methods lowering error to a negligible level and by this way improving the results. Zsako and Zsako have suggested a closed form approximation for the exponential integral . Although the determination of the activation energy is straightforward, one cannot carry out further kinetic computations without further data or assumptions as to the form of the kinetic equation. It has often been observed that the initial portion of a TG curve can usually be well-fitted by a first-order reaction equation. Since the interest in lifetime or stability studies is most often concerned with the initial portion (low conversion level), it has become a procedural practice to assume first-order kinetics and compute the rate constants, pre-exponential factors, etc. on this assumption. When the value of n = 1 is used (first order kinetics), the method is essentially the same as that used by ASTM Standart Test Method E 1641.\n\nIn this work, the mathematical treatment followed in calculating the parameters are called as Standard method. A head start of this method is to obtain the activation energy exactly independently of any assumptions regarding the relation between the reaction rate and the percent reacted. The quality of the statistical fit can be judged from the 95% confidence limits and from observing the lnβ vs 1/T plot. In the ideal case, the data points fall on parallel straight lines, the slope of which is used to calculate the activation energy.\n\n2. Experimental\n\nApparatus\n\nLow heating rate degradation experiments were conducted in a Perkin Elmer PYRIS Diamond TG/DTA (DSC) apparatus in a dynamic nitrogen atmosphere which was calibrated with calcium oxalate. All TGA experiments were carried out with a linear heating rate of 5˚C - 25˚C/min up to 600˚C under nitrogen atmosphere. Kinetic parameters of degradation were determined by the Scanning Kinetics Software.\n\nPolystyrene with Mw = 230,000 g/mol (430,102), poly(p-methylstyrene) with Mw = 72000 (182,273), poly(p-chlorostyrene) with Mw = 75,000 (434,124), and poly(p-bromostyrene) with Mw = 65,000 (181,366) beads were all supplied from Aldrich company and used without any further process. TG analysis was carried out on approximately 10 mg samples.\n\n3. Result and Discussion\n\n3.1. Effect of Heating Rate on p-Substituted PSs\n\nNonisothermal TGA of the samples was carried out at 5, 10, 15, 20, and 25 K/min. Thermogravimetric curves for PS, PMS, PClS, and PBrS are represented in Figure 1 and Figure 2. The kinetic analysis of the degradation of samples was performed by utilizing the thermograms obtained from the TGA. Increase in the heating rate results in a shift toward higher temperatures in the thermograms (Figure 1 and Figure 2). The reason for this behavior is that the sample needs shorter time to catch a given temperature for a faster heating rate.\n\nSame behavior with thermograms was observed for the peak temperatures. An empirical relation between the peak volatilization temperature (Tpv in ˚C) and the heating rate (R in ˚C/min) was reported by Liu et al. .", null, "(8)\n\nThe temperature at which the decomposition rate of the sample is maximum is called as peak volatilization temperature. Tpv values for the PS and the para-substituted derivatives at different heating rates were calculated by using this relationship in the current study, and compared with the TGA experimental data. Experimental values are quite close with that of Liu et al.’s model (Table 1).\n\nln(k) versus 1/T plots for poly(para-substituted styrene)s are shown in Figure 3. The lines obtained at different conversion rates are observed to be roughly parallel to one another and the correlation coefficient of the li-\n\nFigure 1. TGA scans of poly(p-substituted styrene)s at different heating rates in nitrogen.\n\nFigure 2. DTG curves of poly(p-substituted styrene)s at different heating rates in nitrogen.\n\nFigure 3. Kinetic analysis using Ozawa’s method for poly(p-substituted styrene)s degradation in nitrogen environment.\n\nTable 1. Comparison of experimental poly(p-substituted styrene)s peak volatilization decomposition temperatures (Tpv) measured and expected values using Liu’s model in nitrogen at different heating rates.\n\nnearity for the calculation of each activation energy are larger than 0.98. Also there is no systematic divergence from linearity. Therefore it is suggested that the method fitted well for the nonisothermal degradation of poly(p- substituted styrene)s and rate order of degradation is 1 in all poly(p-substituted styrene)s.\n\n3.2. Effect of Substituent Group on Activation Energy and lnZ Values for the Decomposition of p-Substituted PSs\n\nThe activation energies and lnZ values of PS, PMS, PClS, and PBrS at the percent conversion levels of 3.0, 5.0, 8.0, 12.0, 18.0, and 26.0 are tabulated in Table 2. Activation energy is observed to increase with the degree of conversion for PS, PMS, and PClS, which should be attributed to the nature of random scission degradation and the differences in deviation from stationary reaction state at the different heating rates .\n\nThe electronegativity values of methyl and hydrogen which are both bonded to the same group are close to each other. The bonding character of methyl group in PMS is most likely to bonding character of hydrogen in PS. The only difference between the two polymers is the weight of the repeating unit in the backbone. The repeating segment in PMS is slightly heavier in comparison with PS. Therefore the activation energies of PMS are slightly lower than the activation energies of PS. In the case of PClS, the weight of the repeating elements with Cl atom is much higher than the repeating units of PS and PMS and so the activation energy of the PClS is much lower. However PBrS exhibits completely different attitude. Its activation energy is considerably higher than PS and it decreases with the weight loss value. Therefore it can be suggested that the degradation mechanism of PBrS is different.\n\nIn order to understand the dependency of activation energies to the substituent in poly(p-substituted styrene)s, their activation energy values were graphed with respect to molecular weight of the repeating unit in Figure 4.\n\nTable 2. Kinetic analysis results for poly(p-substituted styrene)s.\n\nFigure 4. Activation energies of poly(p-substituted styrene)s at different weight loss percentages.\n\nThe outcome of Figure 4 is as follows;\n\n1) The activation energies of degradation of PS and PMS are keeping up with its values upon different weight loss values as the activation energies of PClS and PBrS are considerably affected by the weight loss ratio.\n\n2) The activation energies of PS, PMS, and PClS increases with the weight loss ratio, as the activation energy value of PBrS decreases.\n\n3) The activation energy of the substituted styrenes decreases with the molecular weight of the repeating units, PBrS shows an extraordinary behavior. Its activation energy values are higher than the others up to high values of the weight losses.\n\nUpon handling individually, PS degrades with chain scission mostly. The bonding character of methyl group in PMS is very similar to bonding character of H in PS. The electronegativity values of methyl and H which are both bonded to the same group are close to each other. The only difference between the two polymers is the weight of the repeating unit in the backbone. The weight of repeating segment of PMS is slightly higher than that of PS. Therefore the activation energies of PMS are slightly lower than PS. In the case of PClS, the weight of the repeating elements with Cl atom is much higher than the repeating units of PS and PMS and the activation energy of the PClS is much lower.\n\nThe bonds between the substituent and the polystyrene backbone are C-H in PS, C-C in PMS, C-Cl in PClS, and C-Br in PBrS. The trend in the electronegativity values of the substituent of polystyrenes is as H(2.20) < C(2.55) < Br(2.96) < Cl(3.16) and, the bonding character of Br is much close to H and C than Cl. As a result, it is seen from the Figure 4 that the activation energy values of the degradation of PBrS is more close to the PS and PMS for any of the weight loss rates. On the other hand the character of Br substituent makes the chemical bonds between repeating groups stronger, leading to a high activation energy values.\n\n4. Conclusion\n\nThermal decomposition kinetics of poly(p-substituted styrene)s was examined in nitrogen atmosphere. Five different heating rates were applied to the samples and as heating rate increased, the degradation shifted to higher temperatures accompanied with an increase in the rate of weight loss. From the nonisothermal degradation data, kinetic parameters were calculated using Flynn-Wall-Ozawa method and the values of the reaction order of poly(p-substituted styrene)s were deduced as 1 with the correlation factor higher than 0.98. Activation energies of thermal degradations are increasing individually in PS, PMS, PClS with weight loss as it is decreasing in PbrS with weight loss. In general, the decrease in the activation energies of poly(p-substituted styrene)s with PS/PMS/ PClS trend is consistent with the increase in repeating unit molecular weight except for PBrS. Therefore PBrS is suggested to decompose with somehow different mechanism.\n\nNOTES\n\n*Corresponding author.\n\nConflicts of Interest\n\nThe authors declare no conflicts of interest.", null, "", null, "", null, "", null, "", null, "[email protected]", null, "+86 18163351462(WhatsApp)", null, "1655362766", null, "", null, "Paper Publishing WeChat", null, "" ]
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https://www.hackmath.net/en/math-problem/10801
[ "# Coal mine\n\nThe towing wheel has a diameter of 1.7 meters. How many meters does the elevator cage lower when the wheel turns 32 times?\n\nCorrect result:\n\nl =  171 m\n\n#### Solution:", null, "We would be pleased if you find an error in the word problem, spelling mistakes, or inaccuracies and send it to us. Thank you!", null, "#### You need to know the following knowledge to solve this word math problem:\n\nWe encourage you to watch this tutorial video on this math problem:\n\n## Next similar math problems:\n\n• Angles 1", null, "It is true neighboring angles have not common arm?\n• Rhombus angles", null, "If one angle in the rhombus is 159°, what is it neighboring angle in rhombus?\n• Circle", null, "What is the radius of the circle whose perimeter is 6 cm?\n• The diameter", null, "The diameter of a circle is 4 feet. What is the circle's circumference?\n• Circle from string", null, "Martin has a long 628 mm string . He makes circle from it. Calculate the radius of the circle.\n• Semicircle", null, "Calculate the length of a semicircle with a radius of 6cm.\n• Round table", null, "A tablecloth should be sewn on a round table with a diameter of 78 cm, which should extend around the table by 10 cm. How many cms of ribbons need to be bought for edging?\n• Bicycle wheel", null, "After driving 157 m bicycle wheel rotates 100 times. What is the radius of the wheel in cm?\n• Circle - simple", null, "The circumference of a circle is 198 mm. How long in mm is its diameter?\n• Circle - easy 2", null, "The circle has a radius 6 cm. Calculate:\n• Clock hands", null, "The second hand has a length of 1.5 cm. How long does the endpoint of this hand travel in one day?\n• Velocipede", null, "The front wheel of velocipede from year 1880 had a diameter 1.8 m. If the front wheel turned again one then rear wheel 6 times. What was the diameter of the rear wheel?\n• Bicycle wheel", null, "Bicycle wheel diameter is 62 cm. How many times turns the bicycle on the road 1 km long?\n• Central angle", null, "What is the length of the arc of a circle with a diameter of 46 cm, which belongs to a central angle of 30°?\n• Well", null, "Rope with a bucket is fixed on the shaft with the wheel. The shaft has a diameter 50 cm. How many meters will drop bucket when the wheels turn 15 times?\n• Athlete", null, "How long length run athlete when the track is circular shape of radius 120 meters and an athlete runs five times in the circuit?\n• Mine", null, "Wheel in traction tower has a diameter 4 m. How many meters will perform an elevator cabin if wheel rotates in the same direction 89 times?" ]
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http://the-dusty-deck.blogspot.com/2011/02/thinking-formally.html
[ "Labels\n\nAI (1) Haskell (1) Java 8 (2) JavaScript (1) Logic (10) Methodology (8) Refactoring (1) Scala (1) Security (2) Tools (9) Xtend (2) Xtext (1)\n\nTuesday, February 15, 2011\n\nThinking Formally\n\nIf there is no alternative to clear thinking about our programs, we should investigate what we can learn from the science of organized thought.  Logic has a long history, with equal roots in philosophy and mathematics.  In some sense, computer programs breath life into abstract logic and give logic the abilitiy to do real things.  And, just like Pinocchio, they often misbehave!\n\nOne of the fundamental concerns of logic is the study of proofs: how can we convince ourselves in a precise and mechanical way of the truth of some claim.  The desire to do so resulted from the need to distinguish between valid and invalid mathematical proofs, expressed using prose.  Sometimes, hidden assumptions were not recognized, leading to faulty proofs.  For example, consider the first axiomatic proof system, Euclidean geomerty, studied in middle school.  It has five axioms, including the famous \"parallel postulate\".  It is an amazing achievement, which is still an excellent example of how to reason formally, and was the source of many important advances in mathematics.  But Euclidean geometry has many hidden assumptions, which are glossed over by the use of informal drawings.  About 2200 years after Euclid, several mathematicians tried to completely formalize geometrical proofs.  David Hilbert's axiomatization of geometry includes twenty axioms instead of Euclid's five.  One of these states that between every two different points on a straight line there is a third point on the same line.  This was so obvious to Euclid that he didn't even think to mention it in his list of axioms.  But axioms are exactly those self-evident propositions on which the theory is based, and the proof isn't logically sound unless all these are stated explicitly.\n\nLogic is closely tied with computation, although the connection was made explicit only towards the 20th century.  A formal proof is a series of deductions, each of which is based on axioms or previously-proven propositions, using a set of deduction rules.  It is supposed to be so simple that it can easily be checked for correctness, without any understanding of its subject matter, whether it is geometry, number theory, or abstract algebra.  In other words, it should be easy to write a program that checks a formal proof; all you need to do is verify that each step follows from the axioms and from previous steps according to the rules.\n\nBut wait!  Is there a limit to the complexity of the axioms or of the rules used in the proof?  Without such limits, the whole exercise falls apart.  For example, imagine a system of axioms that contains all true statements in number theory.  This is, of course, an infinite set.  How difficult is it to write a program to check whether a given statement is an axiom?  Well, it turns out to be impossible!  (This is a consequence of Gödel's first incompleteness theorem.)  So, if we are to be able to write a program to check formal proofs, we need to restrict the possible sets of axioms.  For example, in Gödel's theorems, the set of axioms is required to be recursively enumerable; this means that there is a (finite) program that can generate all axioms and never generates anything that is not an axiom.  Of course, this program will never terminate if the set is infinite.  That's acceptable, as long as the program is guaranteed to produce every axiom if you wait long enough.\n\nSo now, a concept about computability has entered the study of logic itself.  It is only fair, then, to apply logic to computation; and, indeed, the study of logics of programs is an important field in computer science.  There are even computer programs called theorem provers, which have been used to prove difficult mathematical theorems as well as to verify the correctness of other programs.  However, this has still not affected the day-to-day work of most software developers.  In the next post, I will discuss why this is the case and what formal tools can still be used for software development today." ]
[ null ]
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https://www.mecharithm.com/fundamentals-of-robot-motions-configurations-introduction/
[ "Blog Post\n\n# Fundamentals of Robot Motions: Configurations (Introduction)\n\nThis post is part one of the series of lessons on the fundamentals necessary to represent the robot’s configuration, and it gives an introduction to what we mean when we are talking about representing a robot’s configuration.\n\nIn previous lessons, we learned that the robot’s configuration answers the question of where the robot is, and we saw that there are two ways to represent the robot’s configuration: Implicit representation, where the configuration is represented by embedding the curved space in higher-dimensional Euclidean space subject to constraints and explicit representation where configuration is represented with a minimum number of coordinates. You can watch the lesson on configuration and configuration space before preceding this lesson.\n\nWe also advise reading the preliminaries lesson since the lessons are related to each other, and each one complements the other lesson.\n\nFor the full comprehension of the Fundamentals of Robot Motions and the necessary tools to represent the configurations, velocities, and forces causing the motion, please read the following lessons (note that more lessons will be added in the future):\n\nhttps://www.mecharithm.com/category/fundamentals-of-robotics/fundamentals-of-robot-motions/\n\n## Robot’s Configuration on Plane\n\nWe start this lesson with a planar example to see what we mean to express the configuration of a robot. Suppose a toy car with its motion confined to the plane and two coordinate frames {s}, and {b} with their corresponding unit axes (a hat notation shows a unit vector):", null, "A toy car on a plane with two coordinate frames {s} and {b} with the unit axes. The configuration of the toy car can be expressed by finding the position and orientation of the body fixed-frame with respect to the space frame {s}.\n\n{b} is called a body frame since it is a fixed frame attached instantaneously to the moving body. Therefore, to find the configuration of the toy car, we should express the position and orientation of the body’s fixed-frame coordinates in the base frame coordinates.\n\nOne way to represent the orientation of the body coordinates in terms of the base coordinates is using a rotation matrix:\n\n$R = \\begin{pmatrix} \\hat{x}_b.\\hat{x}_s & \\hat{y}_b.\\hat{x}_s \\\\ \\hat{x}_b.\\hat{y}_s & \\hat{y}_b.\\hat{y}_s \\end{pmatrix} = \\begin{pmatrix} cos\\alpha & -sin\\alpha \\\\ sin\\alpha & cos\\alpha \\end{pmatrix}_{2\\times2}$\n\nThe columns of this matrix are the coordinate axes of the {b} frame expressed in the coordinate axes of the {s} frame. The dot represents the dot product between the coordinate axes, and since they are unit vectors, the dot product represents the cosine of the angle between the two vectors.\n\nRecall that if a and b are two vectors with known length and angle between them:", null, "The dot product of two vectors is the multiplication of the length of the two vectors and cosine of the angle between the two vectors.\n\nThen, the dot product of the two vectors can be expressed as the projection of one vector onto the other multiplied by the length of the other vector (in other words, it is the multiplication of the length of the two vectors and cosine of the angle between the two vectors):\n\n$a.b = ||a|| ||b|| cos \\theta$\n\nThe 2 by 2 rotation matrix belongs to the four dimensional space subject to the three constraints as follows:\n\n• Each column of the rotation matrix is a unit vector. In other words, the lengths of the column vectors are 1.\n• Two columns are orthogonal to each other. In other words, the dot product of the columns is zero.\n\nThus, we have one degree of freedom which is parameterized by the angle α to represent the orientation. The rotation matrix is an implicit way to represent the orientation.\n\nNow that we found a way to represent the orientation of the toy car on the plane, it’s time to represent its position. The position of the origin of the body frame in base frame coordinates can be expressed as the vector p from the origin of the {s} frame to the origin of the {b} frame:\n\n$p = p_x \\hat{x}_s + p_y \\hat{y}_s = \\begin{pmatrix} p_x\\\\ p_y \\end{pmatrix} \\in R^2$\n\nThus the configuration of the toy car can be represented by the pair (R,p), which is the description of the orientation and position of {b} with respect to {s}. This means that base frame {s} can be coincident with the body frame {b} by rotating the base frame coordinates around the z-axis by α and then translating the origin by p.\n\n## Robot’s Configuration in Space\n\nUsing a similar approach, we can express the configuration of a robot in space. We can use frames to represent the configuration of a robot in space. One simple way to represent the position and orientation (configuration) of a robot in space is to:\n\n• Fix a frame to the body of the robot according to the Right-Hand Rule (RHR) (sometimes this frame is attached to an important point on the rigid body like the center of mass, but this is not required), and\n• Fix a frame in space", null, "The position and orientation of a robot in space can be represented by the position of the origin of the body frame expressed in the space frame coordinates and the directions of the coordinate axes of the body frame expressed in the space frame coordinates.\n\nThen the robot’s configuration can be represented by the position of the origin of the body frame expressed in the space frame coordinates and the directions of the coordinate axes of the body frame expressed in the space frame coordinates. Note that the body frame is the stationary frame that coincides with the frame attached to the body at a particular instant in time.\n\nThere are several implicit and explicit representations that are used to represent the rotations and configurations in general that we will discuss in the coming lessons. In the next lesson, we will start talking about different ways to represent the orientations in robotics. The next lesson will be dedicated to Rotation matrices which are the implicit representations of orientations.\n\nThe video version of the current lesson can be watched in the link below:\n\nThanks for reading this post. You can also find the other posts on the Fundamentals of Robotics Course in the link below:\n\nhttps://www.mecharithm.com/category/fundamentals-of-robotics/\n\n###### References:\n\n📘 Textbooks:\n\n• Modern Robotics: Mechanics, Planning, and Control by Frank Park and Kevin Lynch\n• A Mathematical Introduction to Robotic Manipulation by Murray, Lee, and Sastry\n\n✍️ Logo design by Minro Art Group\n\n#### If you enjoyed this post, please consider contributing to help us with the website’s running costs and keep making awesome content for you. We deeply thank you for your generous contribution!", null, "Be sure to let us know your thoughts and questions about this post as well as the other posts on the website. You can either contact us through the “Contact” tab on the website or email us at support[at]mecharithm.com." ]
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https://byjus.com/cbse-sample-papers-for-class-9-maths-set-8/
[ "", null, "# CBSE Sample Paper for Class 9 Maths Set 8\n\nMaths is a subject that is purely based on numbers. For some students, it is an interesting subject, but some of them find it difficult. One of the best ways to overcome this Maths fear is to strengthen the basic concepts and practice questions based on them. Once students are done with solving the textbook problems, then they must solve the questions of varying difficulty level. Here we have provided the CBSE Sample Paper Class 9 Maths Set 8 which include different types of Maths problems. Students must solve this paper after completing their syllabus. It will boost their confidence level and help them to score high marks in Maths exam.", null, "", null, "", null, "", null, "", null, "", null, "", null, "", null, "", null, "", null, "", null, "", null, "### Access Other Sets of CBSE Class 9 Maths Sample Papers\n\nStudents can access the different sets of CBSE Class 9 Sample Papers for Maths from the table below:\n\nBy practising the CBSE Class 9 Maths Sample Papers, students can evaluate their current preparation level. Also, they get to know their weak and strong areas.\n\n### Benefits of Solving the CBSE Class 9 Maths Sample Papers\n\nHere we have listed a few benefits of solving the CBSE Sample Papers.\n\n1. Students understand the exam pattern and marking scheme.\n2. By practising the sample papers, students could able to finish the paper on time.\n3. They get to know the types of questions asked in the Maths papers and their difficulty level.\n4. Students can check their current exam preparation by solving the sample papers.\n5. They get confidence to face the exam boldly.\n\nWe hope students find this information on “CBSE Sample Paper for Class 9 Maths Set 8” helpful in exam preparation. Keep learning and stay tuned for further updates on CBSE and other competitive exams. Download BYJU’S App and subscribe to YouTube channel to access interactive Maths and Science videos." ]
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https://docs.ros.org/en/indigo/api/ecl_utilities/html/classecl_1_1UnaryFreeFunction.html
[ "ecl::UnaryFreeFunction< A, R > Class Template Reference\n\nUnary function object for global/static functions. More...\n\n`#include <function_objects.hpp>`\n\nInheritance diagram for ecl::UnaryFreeFunction< A, R >:", null, "[legend]\n\nList of all members.\n\n## Public Member Functions\n\noperator() (A a)\nThis ensures any children objects are deleted correctly.\nUnaryFreeFunction (R(*function)(A))\nUnary function object constructor for global/static functions.\nvirtual ~UnaryFreeFunction ()\n\n## Private Attributes\n\nR(* free_function )(A)\n\n## Detailed Description\n\n### template<typename A, typename R = void> class ecl::UnaryFreeFunction< A, R >\n\nUnary function object for global/static functions.\n\nCreates a function object from a global/static function with a single argument.\n\ngenerateFunctionObject\nTemplate Parameters:\n A : the argument type. R : the return type.\nFunctionObjects.\n\nDefinition at line 200 of file function_objects.hpp.\n\n## Constructor & Destructor Documentation\n\ntemplate<typename A, typename R = void>\n ecl::UnaryFreeFunction< A, R >::UnaryFreeFunction ( R(*)(A) function ) ` [inline]`\n\nUnary function object constructor for global/static functions.\n\nAccepts a global/static function with a single argument and builds the function object around it.\n\nParameters:\n function : a global/static function with a single argument.\n\nDefinition at line 211 of file function_objects.hpp.\n\ntemplate<typename A, typename R = void>\n virtual ecl::UnaryFreeFunction< A, R >::~UnaryFreeFunction ( ) ` [inline, virtual]`\n\nDefinition at line 213 of file function_objects.hpp.\n\n## Member Function Documentation\n\ntemplate<typename A, typename R = void>\n R ecl::UnaryFreeFunction< A, R >::operator() ( A a ) ` [inline, virtual]`\n\nThis ensures any children objects are deleted correctly.\n\nA unary function object call.\n\nRedirects the unary function object call to the composited global/static function.\n\nReturns:\nR : the function's return value.\n\nImplements ecl::UnaryFunction< A, R >.\n\nDefinition at line 222 of file function_objects.hpp.\n\n## Member Data Documentation\n\ntemplate<typename A, typename R = void>\n R(* ecl::UnaryFreeFunction< A, R >::free_function)(A)` [private]`\n\nDefinition at line 225 of file function_objects.hpp.\n\nThe documentation for this class was generated from the following file:\n\necl_utilities\nAuthor(s): Daniel Stonier\nautogenerated on Thu Jun 6 2019 21:17:40" ]
[ null, "https://docs.ros.org/en/indigo/api/ecl_utilities/html/classecl_1_1UnaryFreeFunction__inherit__graph.png", null ]
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https://www.tensorflow.org/hub/tutorials/tf2_image_retraining
[ "Save the date! Google I/O returns May 18-20\n\n# Retraining an Image Classifier\n\n## Introduction\n\nImage classification models have millions of parameters. Training them from scratch requires a lot of labeled training data and a lot of computing power. Transfer learning is a technique that shortcuts much of this by taking a piece of a model that has already been trained on a related task and reusing it in a new model.\n\nThis Colab demonstrates how to build a Keras model for classifying five species of flowers by using a pre-trained TF2 SavedModel from TensorFlow Hub for image feature extraction, trained on the much larger and more general ImageNet dataset. Optionally, the feature extractor can be trained (\"fine-tuned\") alongside the newly added classifier.\n\n### Looking for a tool instead?\n\nThis is a TensorFlow coding tutorial. If you want a tool that just builds the TensorFlow or TF Lite model for, take a look at the make_image_classifier command-line tool that gets installed by the PIP package `tensorflow-hub[make_image_classifier]`, or at this TF Lite colab.\n\n## Setup\n\n``````import itertools\nimport os\n\nimport matplotlib.pylab as plt\nimport numpy as np\n\nimport tensorflow as tf\nimport tensorflow_hub as hub\n\nprint(\"TF version:\", tf.__version__)\nprint(\"Hub version:\", hub.__version__)\nprint(\"GPU is\", \"available\" if tf.test.is_gpu_available() else \"NOT AVAILABLE\")\n``````\n```TF version: 2.4.1\nHub version: 0.11.0\nWARNING:tensorflow:From <ipython-input-1-0831fa394ed3>:12: is_gpu_available (from tensorflow.python.framework.test_util) is deprecated and will be removed in a future version.\nInstructions for updating:\nGPU is available\n```\n\n## Select the TF2 SavedModel module to use\n\nFor starters, use https://tfhub.dev/google/imagenet/mobilenet_v2_100_224/feature_vector/4. The same URL can be used in code to identify the SavedModel and in your browser to show its documentation. (Note that models in TF1 Hub format won't work here.)\n\nYou can find more TF2 models that generate image feature vectors here.\n\nThere are multiple possible models to try. All you need to do is select a different one on the cell below and follow up with the notebook.\n\n``````model_name = \"mobilenet_v3_small_100_224\" # @param ['bit_s-r50x1', 'efficientnet_b0', 'efficientnet_b1', 'efficientnet_b2', 'efficientnet_b3', 'efficientnet_b4', 'efficientnet_b5', 'efficientnet_b6', 'efficientnet_b7', 'inception_v3', 'inception_resnet_v2', 'mobilenet_v2_100_224', 'mobilenet_v2_130_224', 'mobilenet_v2_140_224', 'mobilenet_v3_large_100_224', 'mobilenet_v3_large_075_224', 'mobilenet_v3_small_100_224', 'mobilenet_v3_small_075_224', 'nasnet_large', 'nasnet_mobile', 'pnasnet_large', 'resnet_v1_50', 'resnet_v1_101', 'resnet_v1_152', 'resnet_v2_50', 'resnet_v2_101', 'resnet_v2_152']\n\nmodel_handle_map = {\n\"efficientnet_b0\": \"https://tfhub.dev/tensorflow/efficientnet/b0/feature-vector/1\",\n\"efficientnet_b1\": \"https://tfhub.dev/tensorflow/efficientnet/b1/feature-vector/1\",\n\"efficientnet_b2\": \"https://tfhub.dev/tensorflow/efficientnet/b2/feature-vector/1\",\n\"efficientnet_b3\": \"https://tfhub.dev/tensorflow/efficientnet/b3/feature-vector/1\",\n\"efficientnet_b4\": \"https://tfhub.dev/tensorflow/efficientnet/b4/feature-vector/1\",\n\"efficientnet_b5\": \"https://tfhub.dev/tensorflow/efficientnet/b5/feature-vector/1\",\n\"efficientnet_b6\": \"https://tfhub.dev/tensorflow/efficientnet/b6/feature-vector/1\",\n\"efficientnet_b7\": \"https://tfhub.dev/tensorflow/efficientnet/b7/feature-vector/1\",\n}\n\nmodel_image_size_map = {\n\"efficientnet_b0\": 224,\n\"efficientnet_b1\": 240,\n\"efficientnet_b2\": 260,\n\"efficientnet_b3\": 300,\n\"efficientnet_b4\": 380,\n\"efficientnet_b5\": 456,\n\"efficientnet_b6\": 528,\n\"efficientnet_b7\": 600,\n\"inception_v3\": 299,\n\"inception_resnet_v2\": 299,\n\"nasnet_large\": 331,\n\"pnasnet_large\": 331,\n}\n\nmodel_handle = model_handle_map.get(model_name)\npixels = model_image_size_map.get(model_name, 224)\n\nprint(f\"Selected model: {model_name} : {model_handle}\")\n\nIMAGE_SIZE = (pixels, pixels)\nprint(f\"Input size {IMAGE_SIZE}\")\n\nBATCH_SIZE = 32\n``````\n```Selected model: mobilenet_v3_small_100_224 : https://tfhub.dev/google/imagenet/mobilenet_v3_small_100_224/feature_vector/5\nInput size (224, 224)\n```\n\n## Set up the Flowers dataset\n\nInputs are suitably resized for the selected module. Dataset augmentation (i.e., random distortions of an image each time it is read) improves training, esp. when fine-tuning.\n\n``````data_dir = tf.keras.utils.get_file(\n'flower_photos',\nuntar=True)\n``````\n```Downloading data from https://storage.googleapis.com/download.tensorflow.org/example_images/flower_photos.tgz\n228818944/228813984 [==============================] - 14s 0us/step\n```\n``````datagen_kwargs = dict(rescale=1./255, validation_split=.20)\ndataflow_kwargs = dict(target_size=IMAGE_SIZE, batch_size=BATCH_SIZE,\ninterpolation=\"bilinear\")\n\nvalid_datagen = tf.keras.preprocessing.image.ImageDataGenerator(\n**datagen_kwargs)\nvalid_generator = valid_datagen.flow_from_directory(\ndata_dir, subset=\"validation\", shuffle=False, **dataflow_kwargs)\n\ndo_data_augmentation = False\nif do_data_augmentation:\ntrain_datagen = tf.keras.preprocessing.image.ImageDataGenerator(\nrotation_range=40,\nhorizontal_flip=True,\nwidth_shift_range=0.2, height_shift_range=0.2,\nshear_range=0.2, zoom_range=0.2,\n**datagen_kwargs)\nelse:\ntrain_datagen = valid_datagen\ntrain_generator = train_datagen.flow_from_directory(\ndata_dir, subset=\"training\", shuffle=True, **dataflow_kwargs)\n``````\n```Found 731 images belonging to 5 classes.\nFound 2939 images belonging to 5 classes.\n```\n\n## Defining the model\n\nAll it takes is to put a linear classifier on top of the `feature_extractor_layer` with the Hub module.\n\nFor speed, we start out with a non-trainable `feature_extractor_layer`, but you can also enable fine-tuning for greater accuracy.\n\n``````do_fine_tuning = False\n``````\n``````print(\"Building model with\", model_handle)\nmodel = tf.keras.Sequential([\n# Explicitly define the input shape so the model can be properly\ntf.keras.layers.InputLayer(input_shape=IMAGE_SIZE + (3,)),\nhub.KerasLayer(model_handle, trainable=do_fine_tuning),\ntf.keras.layers.Dropout(rate=0.2),\ntf.keras.layers.Dense(train_generator.num_classes,\nkernel_regularizer=tf.keras.regularizers.l2(0.0001))\n])\nmodel.build((None,)+IMAGE_SIZE+(3,))\nmodel.summary()\n``````\n```Building model with https://tfhub.dev/google/imagenet/mobilenet_v3_small_100_224/feature_vector/5\nModel: \"sequential\"\n_________________________________________________________________\nLayer (type) Output Shape Param #\n=================================================================\nkeras_layer (KerasLayer) (None, 1024) 1529968\n_________________________________________________________________\ndropout (Dropout) (None, 1024) 0\n_________________________________________________________________\ndense (Dense) (None, 5) 5125\n=================================================================\nTotal params: 1,535,093\nTrainable params: 5,125\nNon-trainable params: 1,529,968\n_________________________________________________________________\n```\n\n## Training the model\n\n``````model.compile(\noptimizer=tf.keras.optimizers.SGD(lr=0.005, momentum=0.9),\nloss=tf.keras.losses.CategoricalCrossentropy(from_logits=True, label_smoothing=0.1),\nmetrics=['accuracy'])\n``````\n``````steps_per_epoch = train_generator.samples // train_generator.batch_size\nvalidation_steps = valid_generator.samples // valid_generator.batch_size\nhist = model.fit(\ntrain_generator,\nepochs=5, steps_per_epoch=steps_per_epoch,\nvalidation_data=valid_generator,\nvalidation_steps=validation_steps).history\n``````\n```Epoch 1/5\n91/91 [==============================] - 30s 174ms/step - loss: 1.1899 - accuracy: 0.5899 - val_loss: 0.7321 - val_accuracy: 0.8423\nEpoch 2/5\n91/91 [==============================] - 15s 160ms/step - loss: 0.6630 - accuracy: 0.8935 - val_loss: 0.7036 - val_accuracy: 0.8651\nEpoch 3/5\n91/91 [==============================] - 14s 158ms/step - loss: 0.6405 - accuracy: 0.9095 - val_loss: 0.6973 - val_accuracy: 0.8580\nEpoch 4/5\n91/91 [==============================] - 15s 160ms/step - loss: 0.6143 - accuracy: 0.9156 - val_loss: 0.6817 - val_accuracy: 0.8722\nEpoch 5/5\n91/91 [==============================] - 15s 160ms/step - loss: 0.5917 - accuracy: 0.9323 - val_loss: 0.6795 - val_accuracy: 0.8778\n```\n``````plt.figure()\nplt.ylabel(\"Loss (training and validation)\")\nplt.xlabel(\"Training Steps\")\nplt.ylim([0,2])\nplt.plot(hist[\"loss\"])\nplt.plot(hist[\"val_loss\"])\n\nplt.figure()\nplt.ylabel(\"Accuracy (training and validation)\")\nplt.xlabel(\"Training Steps\")\nplt.ylim([0,1])\nplt.plot(hist[\"accuracy\"])\nplt.plot(hist[\"val_accuracy\"])\n``````\n```[<matplotlib.lines.Line2D at 0x7f0f2030b198>]\n```", null, "", null, "Try out the model on an image from the validation data:\n\n``````def get_class_string_from_index(index):\nfor class_string, class_index in valid_generator.class_indices.items():\nif class_index == index:\nreturn class_string\n\nx, y = next(valid_generator)\nimage = x[0, :, :, :]\ntrue_index = np.argmax(y)\nplt.imshow(image)\nplt.axis('off')\nplt.show()\n\n# Expand the validation image to (1, 224, 224, 3) before predicting the label\nprediction_scores = model.predict(np.expand_dims(image, axis=0))\npredicted_index = np.argmax(prediction_scores)\nprint(\"True label: \" + get_class_string_from_index(true_index))\nprint(\"Predicted label: \" + get_class_string_from_index(predicted_index))\n``````", null, "```True label: daisy\nPredicted label: daisy\n```\n\nFinally, the trained model can be saved for deployment to TF Serving or TF Lite (on mobile) as follows.\n\n``````saved_model_path = f\"/tmp/saved_flowers_model_{model_name}\"\ntf.saved_model.save(model, saved_model_path)\n``````\n```INFO:tensorflow:Assets written to: /tmp/saved_flowers_model_mobilenet_v3_small_100_224/assets\nINFO:tensorflow:Assets written to: /tmp/saved_flowers_model_mobilenet_v3_small_100_224/assets\n```\n\n## Optional: Deployment to TensorFlow Lite\n\nTensorFlow Lite lets you deploy TensorFlow models to mobile and IoT devices. The code below shows how to convert the trained model to TF Lite and apply post-training tools from the TensorFlow Model Optimization Toolkit. Finally, it runs it in the TF Lite Interpreter to examine the resulting quality\n\n• Converting without optimization provides the same results as before (up to roundoff error).\n• Converting with optimization without any data quantizes the model weights to 8 bits, but inference still uses floating-point computation for the neural network activations. This reduces model size almost by a factor of 4 and improves CPU latency on mobile devices.\n• On top, computation of the neural network activations can be quantized to 8-bit integers as well if a small reference dataset is provided to calibrate the quantization range. On a mobile device, this accelerates inference further and makes it possible to run on accelerators like EdgeTPU.\n\n### Optimization settings\n\n```Wrote TFLite model of 6097236 bytes.\n```\n``````interpreter = tf.lite.Interpreter(model_content=lite_model_content)\n# This little helper wraps the TF Lite interpreter as a numpy-to-numpy function.\ndef lite_model(images):\ninterpreter.allocate_tensors()\ninterpreter.set_tensor(interpreter.get_input_details()['index'], images)\ninterpreter.invoke()\nreturn interpreter.get_tensor(interpreter.get_output_details()['index'])\n``````\n``````num_eval_examples = 50\neval_dataset = ((image, label) # TFLite expects batch size 1.\nfor batch in train_generator\nfor (image, label) in zip(*batch))\ncount = 0\ncount_lite_tf_agree = 0\ncount_lite_correct = 0\nfor image, label in eval_dataset:\nprobs_lite = lite_model(image[None, ...])\nprobs_tf = model(image[None, ...]).numpy()\ny_lite = np.argmax(probs_lite)\ny_tf = np.argmax(probs_tf)\ny_true = np.argmax(label)\ncount +=1\nif y_lite == y_tf: count_lite_tf_agree += 1\nif y_lite == y_true: count_lite_correct += 1\nif count >= num_eval_examples: break\nprint(\"TF Lite model agrees with original model on %d of %d examples (%g%%).\" %\n(count_lite_tf_agree, count, 100.0 * count_lite_tf_agree / count))\nprint(\"TF Lite model is accurate on %d of %d examples (%g%%).\" %\n(count_lite_correct, count, 100.0 * count_lite_correct / count))\n``````\n```TF Lite model agrees with original model on 50 of 50 examples (100%).\nTF Lite model is accurate on 47 of 50 examples (94%).\n```\n[{ \"type\": \"thumb-down\", \"id\": \"missingTheInformationINeed\", \"label\":\"Missing the information I need\" },{ \"type\": \"thumb-down\", \"id\": \"tooComplicatedTooManySteps\", \"label\":\"Too complicated / too many steps\" },{ \"type\": \"thumb-down\", \"id\": \"outOfDate\", \"label\":\"Out of date\" },{ \"type\": \"thumb-down\", \"id\": \"samplesCodeIssue\", \"label\":\"Samples / code issue\" },{ \"type\": \"thumb-down\", \"id\": \"otherDown\", \"label\":\"Other\" }]\n[{ \"type\": \"thumb-up\", \"id\": \"easyToUnderstand\", \"label\":\"Easy to understand\" },{ \"type\": \"thumb-up\", \"id\": \"solvedMyProblem\", \"label\":\"Solved my problem\" },{ \"type\": \"thumb-up\", \"id\": \"otherUp\", \"label\":\"Other\" }]" ]
[ null, "https://www.tensorflow.org/hub/tutorials/tf2_image_retraining_files/output_CYOw0fTO1W4x_1.png", null, "https://www.tensorflow.org/hub/tutorials/tf2_image_retraining_files/output_CYOw0fTO1W4x_2.png", null, "https://www.tensorflow.org/hub/tutorials/tf2_image_retraining_files/output_oi1iCNB9K1Ai_0.png", null ]
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https://www.lmfdb.org/EllipticCurve/Q/176400/bq/
[ "# Properties\n\n Label 176400.bq Number of curves $1$ Conductor $176400$ CM no Rank $1$\n\n# Related objects\n\nShow commands for: SageMath\nsage: E = EllipticCurve(\"bq1\")\n\nsage: E.isogeny_class()\n\n## Elliptic curves in class 176400.bq\n\nsage: E.isogeny_class().curves\n\nLMFDB label Cremona label Weierstrass coefficients Torsion structure Modular degree Optimality\n176400.bq1 176400nw1 [0, 0, 0, 105, 12985] [] 147456 $$\\Gamma_0(N)$$-optimal\n\n## Rank\n\nsage: E.rank()\n\nThe elliptic curve 176400.bq1 has rank $$1$$.\n\n## Complex multiplication\n\nThe elliptic curves in class 176400.bq do not have complex multiplication.\n\n## Modular form 176400.2.a.bq\n\nsage: E.q_eigenform(10)\n\n$$q - 5q^{11} + 2q^{13} + 6q^{17} + 2q^{19} + O(q^{20})$$" ]
[ null ]
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https://fconferences.cirm-math.fr/slides-1386.html
[ "GPELab, an open source Matlab toolbox for the numerical simulation of Gross-Pitaevskii equations   (pdf)\n\nQuantized vortex stability and dynamics in superfluidity and superconductivity   (pdf)\n\nClassical and non-classical flows of superfluids   (pdf)\n\nHigh-order numerical schemes for computing the dynamics of nonlinear Schrödinger equation   (pdf)\n\nHelicity, Topology and Kelvin Waves in reconnecting quantum knots   (pdf)\n\n• Yongyong Cai (Purdue University and Beijing Computational Science Research Center)\n\nGround states and dynamics  for SO coupled Bose-Einstein condensation   (pdf)\n\nOn time splitting for NLS in the semiclassical regime   (pdf)\n\nCounterflowing superfluids   (pdf)\n\nAn overview of the BECASIM project: open source numerical simulators for the Gross-Pitaevskii equation   (pdf)\n\nInhomogeneities and temperature effects in Bose-Einstein condensates   (pdf)\n\nLocalization of nonlocal continuum models   (pdf)\n\nThree-dimensional vortex structures in a rotating dipolar Bose-Einstein condensate   (pdf)\n\nSymmetries and dynamics in a quantum-chaotic system   (pdf)\n\nCoherent optical manipulation of plasma waves   (pdf)\n\nNearly parallel vortex filaments in the 3Gross-Pitaevskii equation   (pdf)\n\nQuantum nature and statistical law in quantum turbulence   (pdf)\n\nSuperfluidity and Bogoliubov theory: rigorous results   (pdf)\n\nA hybrid code for solving the Gross-Pitaevskii equation   (pdf)\n\nOn Some Variational Optimization Problems in Classical Fluids and Super fluids   (pdf)\n\nSingle/Multi Component Quantum Gases: Non-Equilibrum models & applications to Experiments   (pdf)\n\nOn the leapfrogging phenomenon in fluid mechanics   (pdf)\n\nNumerical methods on simulating dynamics of the nonlinear Schrödinger equation  with rotation and/or nonlocal interactions   (pdf)\n\nHigh-order commutator-free Magnus integrators for non-autonomous linear evolution equations   (pdf)\n\nInhomogeneous quantum turbulence in a channel   (pdf)\n\nAn efficient splitting Fourier pseudospectral method for Vlasov-Poisson-Fokker-Planck system   (pdf)\n\nFractional Schrödinger equation: stationary states and dynamics   (pdf)" ]
[ null ]
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http://python.omics.wiki/multiprocessing_map/multiprocessing_partial_function_multiple_arguments
[ "### pool.map - multiple arguments\n\npool.map accepts only a list of single parameters as input. Multiple parameters can be passed to pool by a list of parameter-lists, or by setting some parameters constant using partial.\n\n## Example 1:  List of lists\n\nA list of multiple arguments can be passed to a function via pool.map\n(function needs to accept a list as single argument)\n\nExample: calculate the product of each data pair\n\n`import multiprocessing`\n`import numpy as np`\n\n`data_pairs = [ [3,5], [4,3], [7,3], [1,6] ]`\n\n# define what to do with each data pair ( p=[3,5] ), example: calculate product\n`def myfunc(p):`\n`    product_of_list = np.prod(p)`\n`    return product_of_list`\n\n`if __name__ == '__main__':`\n`    pool = multiprocessing.Pool(processes=4)`\n`    result_list = pool.map(myfunc, data_pairs)`\n`    print(result_list)`\n\n[15, 12, 21, 6]\n\n## Example 2: using partial()\n\nParallel run of a function with multiple arguments\n\nTo use pool.map for functions with multiple arguments, partial can be used to set constant values to all arguments which are not changed during parallel processing, such that only the first argument remains for iterating. (The variable input needs to be always the first argument of a function, not second or later arguments).\n\nExample: multiply all numbers in a list by 10\n\n`import multiprocessing`\n`from functools import partial`\n\n`data_list = [1, 2, 3, 4]`\n\n`def prod_xy(x,y):`\n`    return x * y`\n\n`def parallel_runs(data_list):`\n`    pool = multiprocessing.Pool(processes=4)`\n`    prod_x=partial(prod_xy, y=10) `# prod_x has only one argument x (y is fixed to 10)\n`    result_list = pool.map(prod_x, data_list) `\n`    print(result_list)`\n\n`if __name__ == '__main__':`\n`    parallel_runs(data_list)`\n\n[10, 20, 30, 40]\n\nPartial creates a new simplified version of a function with part of the arguments fixed to specific values." ]
[ null ]
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https://www.studysmarter.us/textbooks/math/precalculus-enhanced-with-graphing-utilities-6th/a-preview-of-calculus-the-limit-derivative-and-integral-of-a-function/q-29-graph-each-function-use-the-graph-to-find-the-indicated/
[ "• :00Days\n• :00Hours\n• :00Mins\n• 00Seconds\nA new era for learning is coming soon", null, "Suggested languages for you:\n\nEurope\n\nAnswers without the blur. Sign up and see all textbooks for free!", null, "Q. 29\n\nExpert-verified", null, "Found in: Page 876", null, "### Precalculus Enhanced with Graphing Utilities\n\nBook edition 6th\nAuthor(s) Sullivan\nPages 1200 pages\nISBN 9780321795465", null, "# Graph each function. Use the graph to find the indicated limit, if it exists. $\\underset{x\\to \\frac{\\mathrm{\\pi }}{2}}{\\mathrm{lim}}f\\left(x\\right),f\\left(x\\right)=\\mathrm{sin}\\left(x\\right)$\n\nThe limit is $1$.\n\nSee the step by step solution\n\n## Step 1. Given Information\n\nWe are given a function $f\\left(x\\right)=\\mathrm{sin}\\left(x\\right)$ and we need to graph the function and find $\\underset{x\\to \\frac{\\mathrm{\\pi }}{2}}{\\mathrm{lim}}f\\left(x\\right)$.\n\n## Step 2. Plotting the graph\n\nSo, the graph of the function $f\\left(x\\right)=\\mathrm{sin}\\left(x\\right)$ is", null, "## Step 3. Finding the limit\n\nTo find the value of $\\underset{x\\to \\frac{\\mathrm{\\pi }}{2}}{\\mathrm{lim}}f\\left(x\\right)$ we need to find the value of the function at $x=\\frac{\\mathrm{\\pi }}{2}$.\n\n$f\\left(x\\right)=\\mathrm{sin}\\left(x\\right)\\phantom{\\rule{0ex}{0ex}}⇒f\\left(\\frac{\\mathrm{\\pi }}{2}\\right)=\\mathrm{sin}\\left(\\frac{\\mathrm{\\pi }}{2}\\right)\\phantom{\\rule{0ex}{0ex}}⇒f\\left(\\frac{\\mathrm{\\pi }}{2}\\right)=1\\phantom{\\rule{0ex}{0ex}}\\therefore \\underset{x\\to \\frac{\\mathrm{\\pi }}{2}}{\\mathrm{lim}}f\\left(x\\right)=1$", null, "### Want to see more solutions like these?", null, "" ]
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https://www.circuit-fantasia.com/circuit_stories/building_circuits/ammeter/op-amp_ammeter/op-amp_ammeter.htm
[ "circuit-fantasia > circuit stories > building circuits > op-amp ammeter     short\n\n#", null, "Circuit Stories on the Whiteboard...\n\n## How do We Build an Op-amp Ammeter?\n\n### (How to convert the imperfect ammeter into an almost ideal op-amp ammeter...)\n\nClassic explanations. \"In this op-amp circuit, an imperfect ammeter (meter movement) is connected in the negative feedback loop. As the op-amp inputs draw no current, all the input current flows through the meter coil. The op-amp strives to make the voltage difference between the two inputs zero. Since V+ is grounded, this means that V- also will be zero (at virtual ground). As a result, the op-amp ammeter has zero input resistance...\"\n\nBut yet, what is the basic idea behind this circuit? In order to reveal the truth, let's build the circuit on the whiteboard (convert the imperfect ammeter into an almost ideal op-amp one)...", null, "Imagine we have a real voltage source with a voltage V and an internal resistance Ri.\n\nThen suppose that we have to measure the current Ish = V/Ri that flows if you short the real voltage source.\n\nBuilding op-amp ammeter (go to Step 1)\n\nLow Current Measurements\n\nProblem. For this purpose, we break the circuit and connect an ammeter (e.g. a moving coil current meter). Only, this ammeter is imperfect - it has some internal resistance RA, which affects the current. The voltage drop VRA across the ammeter is harmful as it enervates the excitating voltage V. Now the voltage difference V - VRA determines the current instead the voltage V. As a result, the current decreases.\n\nWhat do we do to solve the problem?\n\nBuilding op-amp ammeter (go to Step 2)\n\nRemedy: removing a disturbance by an antidisturbance. What do we do in real life when an obstacle stands in our way? We remove it by an equivalent useful \"antidisturbance\".\n\nHere, the voltage drop VRA across the ammeter is harmful; so, we have to remove it by an \"antivoltage\" -VRA. In other words, we have to add so much voltage to the input (excitating) voltage source, as much as it loses across the ammeter.\n\nImplementation: removing a voltage by an \"antivoltage\". Following the recipe above, let's connect an additional supplementary battery and adjust its voltage so that VRA = -VRA. As a result, the \"harmful\" voltage VRA disappears and the point A becomes a virtual ground! The real voltage source V-Ri is \"fooled\": it doesn't \"understand\" that there is an ammeter connected; it \"thinks\" that its output is shorted.\n\nBuilding op-amp ammeter (go to Step 3)", null, "Actually, the additional source helps the input source injecting exactly as much voltage as it drops across the ammeter.\n\nNote that the two voltage sources are connected in series, in one and the same direction (+ -, + -) so that their voltages are added.", null, "Now, we have only to replace the \"op-amp man\" with a real op-amp. It doses the voltage of the power supply thus producing a compensating voltage -VRA. Thus, the combination of an op-amp and a steady battery acts as a regulated supplementary battery.\n\nThe op-amp \"observes\" the potential of the point A (the difference between the two voltages) and changes instantly its output voltage so that the point A stays always at zero volts (it acts as a virtual ground). Doing that, the op-amp compensates the \"harmful\" voltage drop across the ammeter by copying and adding it to the voltage of the input source (in other words, it \"helps\" the input source).\n\nBuilding op-amp ammeter (go to Step 4)\n\n\"Cleaning\" the circuit diagram (don't do it, if you are a real teacher!) Finally, we may remove all the \"unnecessary\" components of the picture (voltage bars, current loops, power supplies etc.) and thus we will get the classic circuit diagram of the op-amp ammeter:) It will be so simple, small, easy to memorize but yet... nonunderstandable! Even if we adorn this bare circuit diagram with all kinds of formulas, it will remain still nonunderstandable!\n\nBut this humble circuit diagram is extremely convinient for writing; so, we may put it into the latest recipe book on electronics and serve it to the poor students :).", null, "Here, the basic idea is hidden...\n\nBuilding op-amp ammeter\n\nAnalog electronics 2004 - Class 9\n\nHow I revealed the secret of parallel negative feedback circuits" ]
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https://www.bionicturtle.com/forum/threads/var-and-distribution.10289/
[ "What's new\n\n# VaR and distribution\n\n#### Ali Ehsan Abbas\n\n##### New Member\nHi everyone:\n\nA 95% VaR measure that assumes normal distribution cuts off at 1.65 critical z.\n\nIf an alternative distribution entails a 95% VaR at 1.56, what does that tell us about properties of the distribution?\n\nIs is safe to assume it exhibits thinner tails?\n\n#### David Harper CFA FRM\n\n##### David Harper CFA FRM\nStaff member\nSubscriber\nHi @Ali Ehsan Abbas I love seemingly straightforward questions that are thought-provoking, like this!", null, "My view is that if we only know the alternative distribution's 95% VaR is 1.56, then we know very little and actually we cannot say is has thinner tails. To demonstrate, consider the exponential distribution https://en.wikipedia.org/wiki/Exponential_distribution which has the remarkably convenient feature that its kurtosis is 6.0; i.e., the exponential is very heavy tailed. Given it is also parsimonious, we can solve for the lambda, λ, param that gives us a 95% quantile of 1.56, which is about 1.92. So that, because exponential CDF = 1 - exp(-λ*X), we can observe that 95.0% = 1 - exp(-1.92*1.56) and we've illustrated a heavy-distribution with a 95% quantile of 1.56.\n\nI think the shows the \"problem:\" we don't know the mean and variance of the alternative distribution. Please keep in mind the 1.65 refers to a standard normal distribution; i.e., mean, µ, of zero and variance, σ^2, of one. My illustrative exponential distribution has mean of 0.5 and variance of 0.27. Changes in location and scale make the bigger difference here. Which leads to me to an even simpler example: assume a normal distribution with µ=-0.09 and variance of one. It's 95% VaR is 1.56 because (1.56-0.09)/1 = 1.65; yet is has normal tails by definition. In this way, we could manufacture all sorts of non-standard normal distributions with a 95% quantile of 1.56.\n\nNow, if the alternative distribution has mean of zero and variance of one, then I think we can infer the alternative distribution has so-called light tails (kurtosis < 3 or excess kurtosis < 0) ... I hope that helps!\n\n#### Ali Ehsan Abbas\n\n##### New Member\nThanks @David Harper CFA FRM\n\nWouldn't a distribution with mean 0 and variance of 1 by default be a standard normal distribution, hence 1.65 deviate to 95% VaR quantile?!!\n\nI think the key is for the distribution to be elleyptical! In short, an ellyptical distribution with 1.56 deviate would exhibit thinnner loss tail viz. standard normal distribution! Please correct me if i'm wrong.\n\n#### David Harper CFA FRM\n\n##### David Harper CFA FRM\nStaff member\nSubscriber\nHi @Ali Ehsan Abbas Re: Wouldn't a distribution with mean 0 and variance of 1 by default be a standard normal distribution. Nope", null, "a distribution with non-zero skew and/or non-zero excess kurtosis is non-normal. For example, a normal mixture distribution can mix two normals, each with µ=0 and variance of 1.0, such that the normal mixture distribution that mixes them equally has mean of zero, variance of 1.0, but heavy-tails (as do all variance mixture distributions). But that's just to illustrate that the definition of normal depends on different third (skew) and fourth (kurtosis) moments.\n\nRe: elliptical, I like your thinking!", null, "I'm not sure, I have my doubts only because I *think* elliptical includes fat-tailed distributions (https://en.wikipedia.org/wiki/Elliptical_distribution) so I don't mean to be a stick in the mud, but couldn't you just \"shift left\" a heavy-tailed distribution ? ... The \"problem\" with 1.56 deviate in isolation is that you don't know where the mean is, so you don't know if that's in the tail or in the body, or the extreme tail etc! However, as usual, I could be wrong", null, "Thanks,\n\nappend: I gave this a bit more thought, and I'm going to double-down on my previous response. I think the condition should be that the 1.56 is a standardized deviate. Any non-normal deviate, say given by (D), can be standardized with (D - µ)/σ, but as mentioned, to standardize does not imply the distribution normal or somehow \"made normal.\" That is, a standardized 1.56 ought to have lighter-than-normal tails (aka, platykurtosis). To be real strict, since I can still imagine exceptions even in this case, I'd say the condition is actually: if the standardized deviate is 1.56 and the distribution is unimodal. Thanks!\n\nLast edited:" ]
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https://www.groundai.com/project/spectral-graph-wavelets-for-structural-role-similarity-in-networks/
[ "Spectral Graph Wavelets for Structural Role Similarity in Networks\n\n# Spectral Graph Wavelets for Structural Role Similarity in Networks\n\nClaire Donnat, Marinka Zitnik, David Hallac & Jure Leskovec\nDepartment of Statistics, Department of Computer Science, Department of Electrical Engineering\nStanford University\nStanford, CA 94305, USA\n{cdonnat,hallac}@stanford.edu ,{marinka,jure}@cs.stanford.edu\n###### Abstract\n\nNodes residing in different parts of a graph can have similar structural roles within their local network topology. The identification of such roles provides key insight into the organization of networks and can also be used to inform machine learning on graphs. However, learning structural representations of nodes is a challenging unsupervised-learning task, which typically involves manually specifying and tailoring topological features for each node. Here we develop GraphWave, a method that represents each node’s local network neighborhood via a low-dimensional embedding by leveraging spectral graph wavelet diffusion patterns. We prove that nodes with similar local network neighborhoods will have similar GraphWave embeddings even though these nodes may reside in very different parts of the network. Our method scales linearly with the number of edges and does not require any hand-tailoring of topological features.\n\nWe evaluate performance on both synthetic and real-world datasets, obtaining improvements of up to 71% over state-of-the-art baselines.\n\n\\iclrfinalcopy\\newfloatcommand\n\ncapbtabboxtable[][\\FBwidth]\n\n## 1 Introduction\n\nStructural role discovery in graphs focuses on identifying nodes which have topologically similar local neighborhoods (i.e., similar local structural roles) while residing in potentially distant areas of the network ( 1). Such alternative definition of node similarity is very different than more traditional notions (Perozzi et al., 2014; Grover and Leskovec, 2016; Yang et al., 2016; Monti et al., 2016; Kipf and Welling, 2017; Hamilton et al., 2017a, b; Garcia-Duran and Niepert, 2017), which all assume some notion of “smoothness” over the graph and thus consider nodes residing in close network proximity to be similar.\n\nSuch structural role information about the nodes can be used for a variety of tasks, including as input to machine learning problems, or even to identify key nodes in a system (principal “influencers” in a social network, critical hubs in contagion graphs, etc.).\n\nWhen structural roles of nodes are defined over a discrete space, they correspond to different topologies of local network neighborhoods (e.g., edge of a chain, center of a star, a bridge between two clusters). However, such discrete roles must be pre-defined, requiring domain expertise and manual inspection of the graph structure. A more powerful and robust method for identifying structural similarity involves learning a continuous vector-valued structural signature of each node in an unsupervised way. This motivates a natural definition of structural similarity in terms of closeness of topological signatures: For any , nodes and are defined to be -structurally similar with respect to a given distance if: . Thus, a robust structural similarity metric must introduce both an appropriate signature and an adequate distance metric.\n\nWhile several methods have been proposed for structural role discovery in graphs, existing approaches are extremely sensitive to small perturbations in the topology and typically lack one or more desirable properties. They often require manually hand-labeling topological features (Henderson et al., 2012), rely on non-scalable heuristics (Ribeiro et al., 2017), and/or return a single similarity score instead of a multidimensional structural signature (Jin et al., 2011, 2014).", null, "Figure 1: Nodes a and b have similar local structural roles even though they are distant in the graph. While raw spectral graph wavelet signatures/coefficients Ψ of a and b might be very different, we treat them as probability distributions and show that the coefficient distributions are indeed similar.\n\nHere we address the problem of structure learning on graphs by developing GraphWave. Building upon techniques from graph signal processing (Coifman et al., 2006; Hammond et al., 2011; Shuman et al., 2013), our approach learns a structural embedding for each node based on the diffusion of a spectral graph wavelet centered at that node. Intuitively, each node propagates a unit of energy over the graph and characterizes its neighboring topology based on the response of the network to this probe. In contrast to prior work that characterizes the wavelet diffusion as a function of the wavelet scaling parameter, we study how the wavelet diffuses through the network at a given scale as a function of the initial source node. We prove that the coefficients of this wavelet directly relate to graph topological properties. Hence, these coefficients contain all the necessary information to recover structurally similar nodes, without requiring the hand-labeling of features. However, the wavelets are, by design, localized on the graph. Therefore to compare structural signatures for nodes that are far away from each other, typical graph signal processing methods (using metrics like correlation between wavelets or distance) cannot be used without specifying an exact one-to-one mapping between nodes for every pairwise comparison, a computationally intractable task.\n\nTo overcome this challenge, we propose a novel way of treating the wavelets as probability distributions over the graph. This way the structural information is contained in how the diffusion spreads over the network rather than where it spreads. In order to provide vector-valued signatures which can then be used as input to any machine learning algorithm, we embed these wavelet distributions using the empirical characteristic function (Lukacs, 1970). The advantage of empirical characteristic functions is that they capture all the moments of a given distribution. This allows GraphWave to be robust to small perturbations in the local edge structure, as we prove mathematically. Computational complexity of GraphWave is linear in the number of edges, thus allowing it to scale to large (sparse) networks. Finally, we compare GraphWave to several state-of-the-art baselines on both real and synthetic datasets, obtaining improvements of up to 71% and demonstrating how our approach is a useful tool for characterizing structural signatures in graphs.\n\nSummary of contributions. The main contributions of our paper are as follows:\n\n• We develop a novel use of spectral graph wavelets by treating them as probability distributions and characterizing the distributions using empirical characteristic functions.\n\n• We leverage these insights to develop a scalable method (GraphWave) for learning node embeddings based on structural similarity in graphs.\n\n• We prove that GraphWave accurately recovers structurally similar nodes.\n\nFurther related work.\n\nPrior work on discovering nodes with similar structural roles has typically relied on explicit featurization of nodes. These methods generate an exhaustive listing of each node’s local topological properties (e.g., node degree, number of triangles it participates in, number of -cliques, its PageRank score) before computing node similarities based on such heuristic representations. A notable example of such approaches is RolX (Henderson et al., 2012), which aims to recover a soft-clustering of nodes into a predetermined number of distinct roles using recursive feature extraction (Henderson et al., 2011). Similarly, struc2vec (Ribeiro et al., 2017) uses a heuristic to construct a multilayered graph based on topological metrics and simulates random walks on the graph to capture structural information. In contrast, our approach does not rely on heuristics (we mathematically prove its efficacy) and does not require explicit manual feature engineering or hand-tuning of parameters.\n\nAnother line of related work are graph diffusion kernels (Coifman et al., 2006) which have been utilized for various graph modeling purposes (Kondor and Lafferty, 2002; Chung, 2007; Rustamov and Guibas, 2013; Tremblay et al., 2014). However, to the best of our knowledge, our paper is the first to apply graph diffusion kernels for determining structural roles in graphs. Kernels have been shown to efficiently capture geometrical properties and have been successfully used for shape detection in the image processing community (Sun et al., 2009; Ovsjanikov et al., 2010; Aubry et al., 2011). However, in contrast to shape-matching problems, GraphWave considers these kernels as probability distributions over real-world graphs. This is because the graphs that we consider are highly irregular (as opposed to the Euclidean and manifold graphs). Therefore, traditional wavelet methods, which typically analyze node diffusions across specific nodes that occur in regular and predictable patterns, do not apply. Instead, by treating wavelets as distributions, GraphWave characterizes the shape of the diffusion, rather than the specific nodes where the diffusion occurs. This key insight allows us to uncover structural signatures and to discover structurally similar nodes.\n\n## 2 Learning structural signatures\n\nGiven an undirected connected graph with nodes , edges , an adjacency matrix (binary or weighted), and a degree matrix , we consider the problem of learning, for every node , a structural signature representing ’s position in a continuous multidimensional space of structural roles.\n\nWe frame this as an unsupervised learning problem based on spectral graph wavelets (Hammond et al., 2011) and develop an approach called GraphWave that provides mathematical guarantees on the optimality of learned structural signatures.\n\n### 2.1 Spectral graph wavelets\n\nIn this section, we provide background on a spectral graph wavelet-based model (Hammond et al., 2011; Shuman et al., 2013) that we will use in the rest of the paper.\n\nLet be the eigenvector decomposition of the unnormalized graph Laplacian and let () denote the eigenvalues of .\n\nLet be a filter kernel with scaling parameter . For simplicity, we use the heat kernel throughout this paper, but our results apply to any low-pass filter kernel (Shuman et al., 2016). For now, we assume that is given; we develop a method for selecting an appropriate value of in Appendix C.\n\nGraph signal processing (Hammond et al., 2011; Shuman et al., 2013) defines the spectral graph wavelet associated with as the signal resulting from the modulation in the spectral domain of a Dirac signal centered around node . The spectral graph wavelet is given by an -dimensional vector:\n\n Ψa=U Diag(gs(λ1),…,gs(λN))UTδa, (1)\n\nwhere is the one-hot vector for node . For notational simplicity, we drop the explicit dependency of spectral graph wavelet on . The -th wavelet coefficient of this column vector is thus given by\n\nIn spectral graph wavelets, the kernel modulates the eigenspectrum such that the resulting signal is typically localized on the graph and in the spectral domain (Shuman et al., 2013). Spectral graph wavelets are based on an analogy between temporal frequencies of a signal and the Laplacian’s eigenvalues. Eigenvectors associated with smaller eigenvalues carry slow varying signal, encouraging nodes that are geographically close in the graph to share similar values. In contrast, eigenvectors associated with larger eigenvalues carry faster-varying signal across edges. The low-pass filter kernel can thus be seen as a modulation operator that discounts higher eigenvalues and enforces smoothness in the signal variation on the graph.\n\n### 2.2 GraphWave algorithm\n\nFirst we describe the GraphWave algorithm (Alg. 1) and then analyze it in the next section. For every node , ] GraphWave returns a -dimensional vector representing its structural signature, where nodes with structurally similar local network neighborhoods will have similar signatures.\n\nWe first apply spectral graph wavelets to obtain a diffusion pattern for every node (Line 3), which we gather in a matrix . Here, is a matrix, where -th column vector is the spectral graph wavelet for a heat kernel centered at node . In contrast to prior work that studies wavelet coefficients as a function of the scaling parameter , we study them as a function of the network (i.e., how the coefficients vary across the local network neighborhood around the node ).\n\nIn particular, coefficients in each wavelet are identified with the nodes and represents the amount of energy that node has received from node . As we will later show nodes and with similar network neighborhoods have similar spectral wavelet coefficients (assuming that we know how to solve the “isomorphism” problem and find the explicit one-to-one mapping of the nodes from ’s neighborhood to the nodes of the ’s neighborhood). To resolve the node mapping problem GraphWave treats the wavelet coefficients as a probability distribution and characterizes the distribution via empirical characteristic functions. This is the key insight that makes it possible for GraphWave to learn nodes’ structural signatures via spectral graph wavelets.\n\nMore precisely, we embed spectral graph wavelet coefficient distributions into -dimensional space (Line 4-7) by calculating the characteristic function for each node’s coefficients and sample it at evenly spaced points. The characteristic function of a probability distribution is defined as:  (Lukacs, 1970). The function fully characterizes the distribution of because it captures information about all the moments of probability distribution  (Lukacs, 1970). For a given node and scale , the empirical characteristic function of is defined as:\n\n ϕa(t)=1NN∑m=1eitΨma (2)\n\nFinally, structural signature of node is obtained by sampling the 2-dimensional parametric function (Eq. (2)) at evenly spaced points and concatenating the values:\n\n χa=[Re(ϕa(ti)),Im(ϕa(ti))]t1,⋯td (3)\n\nNote that we sample/evaluate the empirical characteristic function at points and this creates a structural signature of size . This means that the dimensionality of the structural signature is independent of the graph size. Furthermore, nodes from different graphs can be embedded into the same space and their structural roles can be compared across different graphs.\n\nDistance between structural signatures. The final output of GraphWave is a structural signature for each node in the graph. We can explore distances between the signatures through the use of the distance on . The structural distance between nodes and is then defined as: By definition of the characteristic function, this technique amounts to comparing moments of different orders defined on wavelet coefficient distributions.\n\nScaling parameter. The scaling parameter determines the radius of network neighborhood around each node (Tremblay et al. (2014); Hammond et al. (2011)). A small value of determines node signatures based on similarity of nodes’ immediate neighborhoods. In contrast, a larger value of allows the diffusion process to spread farther in the network, resulting in signatures based on neighborhoods with greater radii.\n\nGraphWave can also integrate information across different radii of neighborhoods by jointly considering many different values of . This is achieved by concatenating representations , each associated with a scale , where . We provide a theoretically justified method for finding an appropriate range and in Appendix C. In this multiscale version of GraphWave, the final aggregated structural signature for node is a vector with the following form:\n\nComputational complexity. We use Chebyshev polynomials (Shuman et al., 2011) to compute Line 3 in Algorithm 1. As in Defferrard et al. (2016), each power of the Laplacian has a computational cost of , yielding an overall complexity of where denotes the order Chebyshev polynomial approximation. The overall complexity of GraphWave is linear in the number of edges, which allows GraphWave to scale to large sparse networks.\n\n## 3 Analysis of GraphWave\n\nIn this section, we provide theoretical motivation for our spectral graph wavelet-based model (Shuman et al., 2013). First we analytically show that spectral graph wavelet coefficients characterize the topological structure of local network neighborhoods (Section 3.1). Then we show that structurally equivalent/similar nodes have near-identical/similar signatures (Sections 3.2 and 3.3), thereby providing a mathematical guarantee on the optimality of GraphWave.\n\n### 3.1 Spectral graph wavelets as a measure of network structure\n\nWe start by establishing the relationship between the spectral graph wavelet of a given node and the topological properties of local network neighborhood centered at . In particular, we prove that a wavelet coefficient provides a measure of network connectivity between nodes and .\n\nWe use the fact that the spectrum of the graph Laplacian is discrete and contained in the compact set . It then follows from the Stone-Weierstrass theorem that the restriction of kernel to the interval can be approximated by a polynomial. This polynomial approximation, denoted as , is tight and its error can be uniformly bounded. Formally, this means:\n\n ∀ϵ>0,∃P:P(λ)=K∑k=0αkλk such that |gs(λ)−P(λ)|≤ϵ∀λ∈[0,λmax], (4)\n\nwhere is the order of polynomial approximation, are coefficients of the polynomial, and is the residual. We can now express the spectral graph wavelet for node in terms of the polynomial approximation as:\n\n Ψa=(K∑k=0αkLk)δa+Ur(Λ)UTδa. (5)\n\nWe note that is a function of and thus can be interpreted using graph theory. In particular, it contains terms of the form (capturing the degree), (capturing the number of -length paths that node participates in), and terms containing both and , which denote paths of length up to going from node to every other node .\n\nUsing the Cauchy-Schwartz’s inequality and the facts that is unitary and is uniformly bounded (Eq. (4)), we can bound the second term on the right-hand side of Eq. (5) by:\n\n |δTmUr(Λ)UTδa|2=|N∑j=1r(λj)UajUmj|2≤(N∑j=1|r(λj)|2U2aj)(N∑j=1 U2mj)≤ϵ2. (6)\n\nAs a consequence, each wavelet can be approximated by a -th order polynomial that captures information about the -hop neighborhood of node . The analysis of Eq. (5), where we show that the second term is limited by , indicates that spectral graph wavelets are predominately governed by topological features (specifically, degrees, cycles and paths) according to the specified heat kernel. The wavelets thus contain the information necessary to generate structural signatures of nodes.\n\n### 3.2 Signatures of structurally equivalent nodes\n\nLet us consider nodes and whose -hop neighborhoods are identical (where is an integer less than the diameter of the graph), meaning that nodes and are structurally equivalent. We now show that and have -structurally similar signatures in GraphWave.\n\nFirst, we use the Taylor expansion to obtain an explicit -th order polynomial approximation of as: . Then, for each eigenvalue , we use the Taylor-Lagrange equality to ensure the existence of such that:\n\n |r(λ)|=|e−λs−P(λ,s)|=(λs)K+1(K+1)!e−λcλ≤(λs)K+1(K+1)!. (7)\n\nIf we take any such that it satisfies: , then the absolute residual in Eq. (7) can be bounded by for each eigenvalue . Here, is a parameter that we can specify depending on how close we want the signatures of structurally equivalent nodes to be (note that smaller values of the scale lead to smaller values of and thus tighter bounds).\n\nBecause and are structurally equivalent, there exists a one-to-one mapping from the -hop neighborhood of (i.e., ) to the -hop neighborhood of (i.e., ), such that: . We extend the mapping to the whole graph by randomly mapping the remaining nodes. Following Eq. (5), we write the difference between each pair of mapped coefficients and in terms of the -th order approximation of the graph Laplacian:\n\n |Ψma−Ψπ(m)b|=∣∣δmU(P(Λ)+r(Λ))UTδa−δπ(m)U(P(Λ)+r(Λ))UTδb∣∣≤∣∣(UP(Λ)UT)ma−(UP(Λ)UT)π(m)a∣∣+∣∣(Ur(Λ)UT)ma∣∣+∣∣(Ur(Λ)UT)π(m)b∣∣. (8)\n\nHere, we analyze the first term on the second line in Eq. (8). Since the -hop neighborhoods around and are identical and by the localization properties of the -th power of the Laplacian (-length paths, Section 3.1), the following holds:\n\n ∀m∈NK(a),(K∑k=0αkLk)ma=(K∑k=0αkLk)π(m)b,∀m∉NK(a),(K∑k=0αkLk)ma=(K∑k=0αkLk)π(m)a=0,\n\nmeaning that this term cancels out in Eq. (8). To analyze the second and third terms on the second line of Eq. (8), we use bound for the residual term in the spectral graph wavelet (Eq. (6)) to uniformly bound entries in matrix by . Therefore, each wavelet coefficient in is within of its corresponding wavelet coefficient in , i.e., As a result, because similarity in distributions translates to similarity in the resulting characteristic functions (Lévy’s continuity theorem), then assuming the appropriate selection of scale, structurally equivalent nodes have -structurally similar signatures.\n\n### 3.3 Signatures of structurally similar nodes\n\nWe now analyze structurally similar nodes, or nodes whose -hop neighborhoods are identical up to a small perturbation of the edges. We show that such nodes have similar GraphWave signatures.\n\nLet denote a perturbed -hop neighborhood of node obtained by rewiring edges in the original -hop neighborhood . We denote by the graph Laplacian associated with that perturbation. We next show that when perturbation of a node neighborhood is small, the changes in the wavelet coefficients for that node are small as well.\n\nFormally, assuming a small perturbation of the graph structure (i.e., for all ), we use -th order Taylor expansion of kernel to express the wavelet coefficients in the perturbed graph as:\n\n ˜Ψa=K∑k=0αk~Lk+~Ur(~Λ)~UT. (9)\n\nWe then use the Weyl’s theorem (Coburn et al., 1966) to relate perturbations in the graph structure to the change in the eigenvalues of the graph Laplacian. In particular, a small perturbation of the graph yields small perturbations of the eigenvalues. That is, for each , is close its original value : , where is a constant. Taking everything together, we get:\n\n |Ψma−˜Ψma|≤|K∑k=0αk(Lk−~Lk)ma|+|~Ur(~Λ)~UT|ma+|Ur(Λ)UT|ma=(K∑k=0|αk|+1+C)ϵ,\n\nindicating that structurally similar nodes have similar signatures in GraphWave.\n\n## 4 Experiments on synthetic graphs\n\nBaselines. We compare our GraphWave method against two state-of-the-art baselines, struc2vec (Ribeiro et al., 2017) and RolX (Henderson et al., 2012). We note that RolX requires the number of desired structural classes as input, whereas the two other methods learn embeddings that capture a continuous spectrum of roles rather than discrete classes. We thus use RolX as an oracle estimator, providing it with the correct number of classes111Project website is: http://snap.stanford.edu/graphwave.. We also note that homophily-based methods (Kipf and Welling (2017); Hamilton et al. (2017a), etc.) are unable to recover structural similarities.\n\n### 4.1 Barbell graph\n\nWe consider a barbell graph consisting of two dense cliques connected by a long chain (Figure 2A). We run GraphWave, RolX, and struc2vec and plot a 2D PCA representation of learned structural signatures in Figure 2B-D.\n\nGraphWave correctly learns identical representations for structurally equivalent nodes, providing empirical evidence for our theoretical result in Section 3.2. This can be seen by structurally equivalent nodes in Figure 2A (nodes of the same color) having identical projections in the PCA plot (Figure 2D). In contrast, both RolX and struc2vec fail to recover the exact structural equivalences.\n\nAll three methods correctly group the clique nodes (purple) together. However, only GraphWave correctly differentiates between nodes connecting the two dense cliques in the barbell graph, providing empirical evidence for our theoretical result in Section 3.3. GraphWave represents those nodes in a gradient-like pattern that captures the spectrum of structural roles of those nodes (Figure 2D).", null, "Figure 2: Barbell graph. The graph has 8 distinct classes of structurally equivalent nodes as indicated by color (A). 2D PCA projection of structural signatures as learned by RolX (B), struc2vec (C) and GraphWave (D). Projections in B-D contain the same number of points as there are nodes in the barbell graph. Identical signatures have identical projections, resulting in overlapping points in B-D.\n\n### 4.2 Complex and noisy graph structures\n\nGraphs. We next consider four types of synthetic graphs where the structural role of each node is known and used as ground truth information to evaluate performance. The graphs are given by basic shapes of one of different types (“house”, “fan”, “star”) that are regularly placed along a cycle (Table 1 and Figure 3A). In the “varied” setup, we mix the three basic shapes when placing them along a cycle, thus generating synthetic graphs with richer and more complex structural role patterns. Additional graphs are generated by placing these shapes irregularly along the cycle followed by adding a number of edges uniformly at random. In our experiments, we set this number to be around % of the edges in the original structure. This setup is designed to assess the robustness of the methods to data perturbations (“house perturbed”, “varied perturbed”).\n\nExperimental setup. For each graph, we run RolX, struc2vec, and GraphWave to learn the signatures. We choose to use a multiscale version of GraphWave where the scale was set as explained in Appendix C. We then use -means to cluster the learned signatures and use three standard metrics to evaluate the clustering quality. (1) Cluster homogeneity is the conditional entropy of the ground-truth structural roles given the proposed clustering (Rosenberg and Hirschberg, 2007). (2) Cluster completeness (Rosenberg and Hirschberg, 2007) evaluates whether nodes with the same structural role are in the same cluster. (3) Silhouette score compares the mean intra-cluster distance to the mean between-cluster distance, assessing the density of the recovered clusters. This score takes a value in [-1,1] (higher is better).\n\nResults. GraphWave consistently outperforms struc2vec, yielding improvements for the homogeneity of up to 50%, and completeness up to 69% in the “varied” setting (Table 1). Both GraphWave and RolX achieved perfect performance in the noise-free “house” setting, however, GraphWave outperformed RolX by up to 4% (completeness) in the more complex “varied” setting. We evaluated methods on graphs in the presence of noise (“perturbed” in Table 1): GraphWave outperformed RolX and struc2vec by 10% and 67% (completeness), respectively, providing empirical evidence for our analytical result that GraphWave is robust to noise in the edge structure. The silhouette scores also show that the clusters recovered by GraphWave are denser and better separated than for the other methods.", null, "Figure 3: A cycle graph with attached “house” shapes (A). 2D PCA projection of GraphWave’s structural signatures. The representation of structurally equivalent nodes overlap, and GraphWave perfectly recovers the 6 different node types (B). Empirical characteristic function for the distribution of the wavelet coefficients (C). Color of a node/curve indicates structural role. (Best seen in color.)", null, "Table 1: Structural role discovery results for different synthetic graphs. (Best seen in color.) Results averaged over 20 synthetically generated graphs. Dashed lines denote perturbed versions of the basic shapes (obtained by randomly adding and removing edges), node colors indicate structural roles.\n\nAs an example, we show a cycle graph with attached “house” shapes (Figure 3A). We plot 2D PCA projections of GraphWave’s signatures in Figure 3B, confirming that GraphWave accurately distinguishes between nodes with distinct structural roles. We also visualize the resulting characteristic functions (Eq. (2)) in Figure 3C. In general, their interpretation is as follows (Appendix D):\n\n• Nodes located in the periphery of the graph struggle to diffuse the signal over the graph, and thus span wavelets that are characterized by a smaller number of non-zero coefficients. Characteristic functions of such nodes thus span a small loop-like 2D curve.\n\n• Nodes located in the core (dense region) of the graph tend to diffuse the signal farther away and reach farther nodes for the same value of . Characteristic functions of such nodes thus have a farther projection on the and axis.\n\nIn Figure 3C, different shapes of the characteristic functions capture different structural roles. We note the visual proximity between the roles of the blue, light green and red nodes that these curves carry, as well as their clear difference with the core dark green and purple nodes.\n\n## 5 Experiments on two real-world graphs\n\n### 5.1 The Enron email graph\n\nData and setup. Nodes represent Enron employees and edges correspond to email communication between the employees (Klimt and Yang, 2004). An employee has one of seven functions in the company (e.g., CEO, president, manager). These functions provide ground-truth information about roles of the corresponding nodes in the network. We use GraphWave to learn a structural signature for every Enron employee. We then use these signatures to compute the average distance between every two categories of employees.", null, "Figure 4: Heat maps indicate average distance between roles in the Enron email graph, as determined by struc2vec (A) and GraphWave (B). Table of comparative statistics for the three algorithms.", null, "Figure 5: PCA projection of the learned airport structural signatures. A: Air France, a major airline, and B: Ryanair, a low-cost fare airline. Selected nodes are labeled using three-letter airports codes.\n\nResults. GraphWave captures intricate organizational structure of Enron (Figure 4). For example, CEOs and presidents are structurally distant from all other job titles. This indicates their unique position in the email exchange graph, which can be explained by their local graph connectivity patterns standing out from the others. Traders, on the other hand, appear very far from presidents and are closer to directors. In contrast, struc2vec is less successful at revealing intricate relationships between the job titles, yielding an almost uniform distribution of distances between every class.\n\nWe assess the separation between “top” job titles (CEO and President) and lower levels in the job title hierarchy. GraphWave achieves 28% better homogeneity and 139% better completeness than RolX. We also note that the variability within each cluster of struc2vec is higher than the average distance between clusters (dark green colors on the diagonal, and lighter colors on the off-diagonal).\n\n### 5.2 European airline graphs\n\nData and setup. The airline graphs are taken from the list of airlines operating flights between European airports (Cardillo et al., 2013). Each airline is represented with a graph, where nodes represent airports and links stand for direct flights between two airports. Given an airline graph, we use GraphWave to learn structural signature of every airport. We then create a visualization using PCA on the learned signatures to layout the graph on a two-dimensional structural space.\n\nResults. Figure 5 shows graph visualizations of two very different airlines. Air France is a national French airline, whose graph has the so-called hub and spoke structure, because the airline is designed to provide an almost complete coverage of the airports in France. We note that the signature of CDG (Charles De Gaulle, which is Air France’s central airport) clearly stands out in this 2D projection, indicating its unique role in the network. In contrast to Air France, Ryanair is a low-cost airline whose graph avoids the overly centralized structure. Ryanair’s network has near-continuous spectrum of structural roles that range from regional French airports (Lille (LIL), Brest Bretagne (BES)) all the way to London Stansted (STN) and Dublin (DUB). These airlines have thus developed according to different structural and commercial constraints, which is clearly reflected in their visualizations.\n\n## 6 Conclusion\n\nWe have developed a new method for learning structural signatures in graphs. Our approach, GraphWave, uses spectral graph wavelets to generate a structural embedding for each node, which we accomplish by treating the wavelets as a distributions and evaluating the resulting characteristic functions. Considering the wavelets as distributions instead of vectors is a key insight needed to capture structural similarity in graphs.\n\nOur method provides mathematical guarantees on the optimality of learned structural signatures. Using spectral graph theory, we prove that structurally equivalent/similar nodes have near-identical/similar structural signatures in GraphWave. Experiments on real and synthetic networks provide empirical evidence for our analytical results and yield large gains in performance over state-of-the-art baselines. For future work, these signatures could be used for transfer learning, leveraging data from a well-explored region of the graph to infer knowledge about less-explored regions.\n\n## References\n\n• Aubry et al. 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In ICML, volume 33, pages 40–48, 2016.\n\n## Notation\n\nIn the appendix, we use the same notation as in the main paper. However, in the main paper, we dropped the explicit dependency of the heat kernel wavelet on the scaling parameter . We note that we explicitly keep this dependency throughout the appendix as we study the relationship between heat kernel wavelets and the scaling parameter.\n\n## Appendix A Known properties of heat kernel wavelets\n\nWe here list five known properties of heat kernel wavelets that we later use to derive a method for automatic selection of the scaling parameter in Appendix C.\n\n1. The heat kernel wavelet matrix is symmetric:\n\n Ψ(s)ma=δTmUe−sΛUTδa=δTaUe−sΛUTδm=Ψ(s)am.\n2. First eigenvector: By definition of the graph Laplacian , the vector is an eigenvector of , corresponding to the smallest eigenvalue . This means that:\n\n U1=1/√N\\mathbbm1––––––––––––⟹UT\\mathbbm1––––––––––––=(√N,0,0,…,0)T=√Nδ1, (**)\n\nwhere the last equality follows from the fact that eigenvectors in are orthogonal.\n\n3. Scaling inequality: Using the Cauchy Schwartz inequality, we get:\n\n (Ψ(s)ma)2=(N∑j=1e−λjsUajUmj)2≤(N∑j=1(e−λjs2Uaj)2)(N∑j=1(e−λjs2Ujm)2)=Ψ(s)aaΨ(s)mm.\n\nIt thus follows that:\n\n4. Convergence: We use the fact that when , the value of converges to 0 for any non-zero eigenvalue of the graph Laplacian. Combining this fact with the spectral graph wavelet definition in Eq. (1), we get:\n\n ∀a,m,lims→∞Ψ(s)am=N∑j=1lims→∞e−sλjUajUmj=Ua1Um1=1/N,\n\nsince it follows from Eq. (**A) that for every . The diffusion over the network thus converges to a state where all nodes have an identical temperature .\n\n5. Convergence rate: To analyze the convergence rate of heat kernel wavelets as a function of , we expand using the spectral graph wavelet definition in Eq. (1):\n\n |Ψ(s+1)aa−1N|=|N∑j=2e−λj(s+1)U2ja|≤|N∑j=2e−λjsU2ja|≤|Ψ(s)aa−1N|,\n\nwhere:\n\n |Ψ(s)aa−1N|=e−λ2s|N∑j=2e−(λj−λ2)sU2ja|. decreasing function of s\n\nTaking everything together, we get: It thus follows that:\n\n## Appendix B Properties of heat kernel wavelet distributions\n\nWe prove three propositions about distributions generated by heat kernel wavelets that will be used in Appendix C.\n\n###### Proposition 1.\n\nThe mean of heat kernel wavelet is equal to . The mean is thus independent of the value of scaling parameter and node .\n\n###### Proof.\n\nWe expand the mean of heat kernel wavelet using spectral graph wavelet definition in Eq. (1) and Property 2 (P2) from Appendix A:\n\n μ(s)a=1NN∑m=1Ψ(s)ma=1N\\mathbbm1––––––––––––TUe−ΛsUδa=1√NδT1e−ΛsUδa=1√NδT1(e−λ1sU1a,e−λ2sU2a,…,e−λNsUNa)T=1√Ne−λ1sU1a=1N.\n\n###### Proposition 2.\n\nThe heat kernel wavelet coefficient at the initial source node is a monotonically decreasing function of scaling parameter . Its value is bounded by: .\n\n###### Proof.\n\nThis follows directly using definition and the scaling inequality in Property 3 (P3). This way, we get:\n\n Ψ(∞)aa=1N≤Ψ(s)aa≤1=Ψ(0)aa.\n\nAdditionally, for any , the wavelet coefficient is non-negative and bounded. Specifically, the wavelet coefficient can be written as: . It can be bounded by: (Property 4).\n\n###### Proposition 3.\n\nThe variance of heat kernel wavelet is a strictly decreasing function of scaling parameter .\n\n###### Proof.\n\nWe use the definition of variance to get:\n\n Var[Ψ(s)a]=1NN∑m=1(Ψ(s)ma−μ(s)a)2=1NN∑m=1(Ψ(s)ma)2−1N2.\n\nWe rewrite the sum on the far right hand-side using the symmetry of wavelet matrix (Property 1):\n\n N∑m=1(Ψ(s)ma)2=N∑m=1Ψ(s)amΨ(s)ma=(Ψ2)aa=δTaUe−2sΛUTδa=||e−sΛUTδa||22=N∑j=1e−2sλjU2aj,\n\nconcluding that variance is decreasing, since it is sum of functions, all of which decrease as gets larger. ∎\n\n## Appendix C Scaling parameter of heat kernel\n\nWe here develop a method that automatically finds an appropriate range of values for the scaling parameter in heat kernel , which we use in the multiscale version of GraphWave (Section 2.2)\n\nWe find the appropriate range of values for by specifying an interval bounded by and through the analysis of variance in heat kernel wavelets. Intuitively, whether or not a given value for is appropriate for structural signature learning depends on the relationship between the scaling parameter and the temporal aspects of heat equation. In particular, small values of allow little time for the heat to propagate, yielding diffusion distributions (i.e., heat kernel wavelet distributions) that are trivial in the sense that only a few coefficients have non-zero values and are thus unfit for comparison. For larger values of the network converges to a state in which all nodes have an identical temperature equal to (Property 4), meaning that diffusion distributions are data-independent, hence non-informative.\n\nNext we prove Propositions 4 and 5 to provide new insights into the variance and convergence rate of heat kernel wavelets. We then use these results to select and .\n\n###### Proposition 4.\n\nGiven the scaling parameter , the variance of off-diagonal coefficients in heat kernel wavelet is proportional to:\n\n Var[{Ψ(s)am;m≠a}]∝Δ(0)aΔ(2s)a−(Δ(s)a)2,\n\nwhere is a monotonically decreasing function of .\n\n###### Proof.\n\nLet us denote the mean of off-diagonal coefficients in wavelet by: . We use the fact that , along with the definition of the variance, to obtain:\n\n Var[{Ψ(s)am;m≠a}]=1N−1∑m≠a(Ψ(s)ma−~μ(s)a)2=1N−1∑m≠a(Ψ(s)ma)2−(~μ(s)a)2=1N−1(Ψ(2s)aa−(Ψ(s)aa)2)−1(N−1)2(1−Ψ(s)aa)2=N(N−1)2(Ψ(2s)aaN−1N−(Ψ(s)aa)2+2Ψ(s)aaN−1N)=N(N−1)2((Ψ(2s)aa−1N)N−1N−(Ψ(s)aa−1N)2)=N(N−1)2(Δ(0)aΔ(2s)a−(Δ(s)a)2).\n\nProposition 4 proves that the variance is a function of . Therefore, to maximize the variance, we must analyze the behavior of . To ensure sufficient variability in the distribution of wavelet coefficients, we need to select a range that bounds the . Our goal thus becomes establishing that is large enough that the diffusion has had time to spread, while remaining sufficiently small to ensure that the diffusion is far from its converged state.\n\n###### Proposition 5.\n\nThe convergence of heat kernel wavelet coefficient is bounded by:\n\n e−λN⌈s⌉Δ(0)a≤Δ(s)a≤e−λ2⌊s⌋Δ(0)a.\n###### Proof.\n\nWe use Property 5 from Appendix A and induction over to complete this proof. For a given we analyze:\n\n |Ψ(s+1)aa−1N|=|N∑j=2e−λj(s+1)U2ja|≤e−λ2|Ψ(s)aa−1N|,\n\nand conclude that: .\n\nGiven any , we use the induction principle to get:\n\n e−λNs|Ψ(0)aa−1N|≤|Ψ(s)aa−1N|≤e−λ2s|Ψ(0)aa−1N|,\n\nwhich immediately yields the desired bound: Since is smooth increasing function of , we can take the floor/ceiling of any non-integer and this proposition must hold. ∎\n\n### c.1 Selection of smax\n\nWe select such that wavelet coefficients are localized in the network. To do so, we use Proposition 5 and bound by the graph Laplacian’s eigenvalues. When the bulk of the eigenvalues leans towards , is closer to (i.e., lower bound in Proposition 5). When the bulk of the eigenvalues is closer to , will lean towards (i.e., upper bound in Proposition 5). In each case, the diffusion is localized if is above a given threshold Indeed, this ensures that has shrunk to at most % of its initial value at , and yields a bound of the form: . The bound implies that: , or\n\nTo find a middle ground between the two convergence scenarios, we take to be the geometric mean of and . Indeed, as opposed to the arithmetic mean, the geometric mean maintains an equal weighting across the range [], and a change of in has the same effect as a change in of . We thus select as:\n\n### c.2 Selection of smin\n\nWe select to ensure the adequate diffusion resolution. In particular, we select a minimum value such that each wavelet has sufficient time to spread. That is, . As in the case of above, we obtain a bound of . Hence, we set to:\n\nTo cover an appropriate range of scales, we suggest setting and .\n\n## Appendix D Visualization of characteristic functions\n\nHere, we study the properties of characteristic functions in GraphWave. Our goal is to provide intuition to understand the behavior of these functions and how their resulting 2D parametric curves reflect nodes’ local topology (see Figure 3C).\n\nWe begin by reviewing the definition of the characteristic function.\n\n###### Definition 1.\n\nThe empirical characteristic function [Lukacs, 1970] of wavelet is function defined as:\n\n ϕ(s)a(t)=1NN∑m=1eitΨ(s)ma,t∈R.\n\nIn the phase plot, a given value of thus yields the following set of coordinates:\n\n ϕ(s)a(t)=⎛⎝1N∑Nm=1cos(tΨ(s)ma)1N∑Nm=1sin(tΨ(s)ma)⎞⎠.\n\nBy varying , we get a characteric curve in this 2D plane. Here, we note several properties of this curve:\n\n1. Value at : , independent of the scale or initial source node.\n\n2. Behavior for : Only the coefficient corresponding to the initial source node has non-zero value (i.e., , and ). Hence, . The phase plot of the characteristic curve is a circle of radius with period centered at .\n\n3. Behavior as : From Property 4 in Appendix A, the coefficients all converge to the same limit:\n\nThis in turn implies that: . Hence, the curve converges to a circle centered at 0, with radius 1 and period .\n\n4. Gradient of the curve: , or equivalently, in the phase plot:\n\n ∇tϕ(s)a(t)=⎛⎝−1N∑Nm=1Ψ(s)masin(tΨ(s)ma)1N∑Nm=1Ψ(s)ma cos(tΨ(s)ma)⎞⎠.\n\nSince all the wavelet coefficients are non-negative, the curve is thus directed counter-clockwise.\n\nYou are adding the first comment!\nHow to quickly get a good reply:\n• Give credit where it’s due by listing out the positive aspects of a paper before getting into which changes should be made.\n• Be specific in your critique, and provide supporting evidence with appropriate references to substantiate general statements.\n• Your comment should inspire ideas to flow and help the author improves the paper.\n\nThe better we are at sharing our knowledge with each other, the faster we move forward.\nThe feedback must be of minimum 40 characters and the title a minimum of 5 characters", null, "", null, "", null, "" ]
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https://essayquip.com/2021/07/19/generate-the-data-from-p-1-1-1-1-1-1-1-1-1-1-x-p-p-p-p-p-p-p-p-for-each-case-what-components-do-you-see-in-the-data-essaylink/
[ "# Generate the data from p = [1 1 1 1 1 -1 -1 -1 -1 -1]; x = [p p p p p p p p]; For each case, what components do you see in the data? – Essaylink\n\nDetermine the 64 point Haar transform for the following sets of data, using the function fasthaar. Then, from the Haar transform data, generate a wavelet map, using the function waveletmap. The data is sample over a period of 5.12 seconds.\n\na. Generate the data from p = [1 1 1 1 -1 -1 -1 -1]; x = [p p p p p p p p]; x = x+0.5*ones(size(x))+0.01*randn(size(x));\n\nb. Generate the data from p = [1 1 1 1 -1 -1 -1 -1]; p1 = [1 1 -1 -1]; x = [p p p1 p1 p1 p1 p p p p];\n\nc. Generate the data from p = [1 1 1 1 1 -1 -1 -1 -1 -1]; x = [p p p p p p p p];\n\nFor each case, what components do you see in the data?", null, "The post Generate the data from p = [1 1 1 1 1 -1 -1 -1 -1 -1]; x = [p p p p p p p p]; For each case, what components do you see in the data? appeared first on Best Custom Essay Writing Services | EssayBureau.com.\n\n0 replies" ]
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https://www.analog.com/en/design-center/reference-designs/circuits-from-the-lab/cn0189.html
[ "Evaluation Hardware \\$60.00", null, "# Overview\n\n## Design Resources\n\n### Design & Integration Files\n\n• Schematic\n• Bill of Materials\n• Gerber Files\n• Assembly Drawing\n\n### Evaluation Hardware\n\nPart Numbers with \"Z\" indicate RoHS Compliance. Boards checked are needed to evaluate this circuit\n\n• EVAL-CFTL-6V-PWRZ (\\$17.00) Wall Power Supply for Eval Board\n• EVAL-CN0189-SDPZ (\\$60.00) Tilt Measurement Using a Dual Axis Accelerometer\n• EVAL-SDP-CB1Z (\\$99.00) Eval Control Board\n\n### Device Drivers\n\nSoftware such as C code and/or FPGA code, used to communicate with component's digital interface.\n\n## Features & Benefits\n\n• 1 Degree Accuracy over 90 Degrees and Temperature\n• Dual axis to ensure precise accuracy\n• Filtered to reduce inband noise\n• Precision tilt applications\n\n# Circuit Function & Benefits\n\nThe circuit, shown in Figure 1, incorporates a dual axis ADXL203 accelerometer and the AD7887 12-bit successive approximation (SAR) ADC to create a dual axis tilt measurement system.\n\nThe ADXL203 is a polysilicon surface micromachined sensor and signal conditioning circuit. Acceleration in the X or Y axis will produce a corresponding output voltage on the XOUT or YOUT output pins of the device. The X axis and Y axis are perpendicular to one another. The AD8608 quad op amp buffers, attenuates, and level shifts the ADXL203 outputs so they are at the proper levels to drive the inputs of the AD7887. The rail-to-rail input/output AD8608 is chosen for its low offset voltage (65 μV maximum), low bias current (1 pA maximum), low noise (8 nV/√Hz), and small footprint (14-lead SOIC or TSSOP).\n\nThe AD7887 is configurable for either dual or single channel operation via the on-chip control register. In this application it is configured for dual channel mode, allowing the user to monitor both outputs of the ADXL203, thereby providing a more accurate and complete solution.\n\nThe system maintains an accuracy of 1° over 90° and over temperature. The circuit provides this precision, performance, and range in a low cost, low power, small footprint, calibration dependent solution. The ADXL203 is specified over a minimum temperature range of −40°C to +105°C and is available in an 8-terminal ceramic leadless chip carrier package (LCC).", null, "Figure 1. Dual Axis Tilt Measurement System (Simplified Schematic: Decoupling and All Connection Not Shown)\n\n# Circuit Description\n\nSupply Voltage and Decoupling\n\nThe ADXL203 requires only one 0.1 μF decoupling capacitor as long as there is no noise present at the 140 kHz internal clock frequency. If necessary, larger bulk capacitors (1 μF to 10 μF) or ferrite beads can be included.\n\nIn order to have output logic levels compatible with the SDP board, the AD7887 must run on a +3.3 V rail. The rest of the circuit, as indicated in Figure 1, uses the +5 V rail. The ADXL203 is specified and tested with a nominal supply voltage of +5 V. Although the ADXL203 is operational with a supply voltage anywhere between 3 V and 6 V, optimum overall performance is achieved at 5 V. Please refer to the ADXL203 data sheet for details regarding performance at other supply voltages.\n\nThe ADXL203 outputs are ratiometric; increasing the supply voltage will act to increase the output voltage. The output sensitivity varies proportionally to supply voltage. At VS = 3 V, the output sensitivity is typically 560 mV/g. At Vs = 5 V, the device has a nominal sensitivity of 1000 mV/g.\n\nThe zero-g output level is also ratiometric, so the zero-g output is nominally equal to VS/2 at all supply voltages.\n\nThe output noise of the ADXL203, however, is not ratiometric but absolute in volts. This means the noise density will decrease as the supply voltage increases. This is because the scale factor (mV/g) increases while the noise voltage remains the same. For VS = 3 V, the noise density is typically 190 μg/√Hz and 110 μg/√Hz for VS = 5 V.\n\nNoise, Bandwidth, and Output Capacitor Selection\n\nThe ADXL203 noise has the characteristics of white Gaussian noise, which contributes equally at all frequencies. It is described in terms of μg/√Hz (the noise is proportional to the square root of the accelerometer bandwidth). The user should limit bandwidth to the lowest frequency needed by the application to maximize the resolution and dynamic range of the accelerometer.\n\nThe bandwidth is set by a capacitor (CX,Y) on the XOUT and YOUT pins of the device. These capacitors create a low-pass filter when combined with the internal 32 kΩ output resistor of the ADXL203. These filters are intended primarily for noise reduction and antialiasing. The equation for the 3 dB bandwidth is:\n\nBW = 1/(2πR×C(X,Y)), where R = 32 kΩ\n\nWith the single poll roll-off characteristic, the typical noise of the ADXL203 on a 5 V supply is determined by:\n\nRMS Noise = (110 μg/√Hz) × √(BW × 1.57)\n\nOften, the peak-to-peak noise is desired as it gives the best estimate of the uncertainty in a single measurement; peak-topeak noise is estimated by multiplying the rms value by 6.\n\nTable 1 gives the bandwidth, rms noise, and peak-to-peak noise for a given filter capacitor. For this circuit, two 10 μF capacitors create a bandwidth of 0.5 Hz. A minimum capacitance of 2000 pF is required in all cases.\n\n Bandwidth (Hz) CxCy (µF) RMS Noise (mg) Peak-to-Peak Noise Estimate (mg) 10 0.47 0.4 2.6 50 0.1 1.0 6 100 0.047 1.4 8.4 500 0.01 3.1 18.7\n\nPhysical Operation of Sensor\n\nThe sensor is a surface micromachined polysilicon structure built on top of the silicon wafer. Polysilicon springs suspend the structure over the surface of the wafer and provide a resistance against acceleration forces. Deflection of the structure is measured using a differential capacitor that consists of independent fixed plates and plates attached to the moving mass.\n\nThe fixed plates are driven by 180° out-of-phase square waves. Acceleration deflects the beam and unbalances the differential capacitor, resulting in an output square wave whose amplitude is proportional to acceleration. Phase-sensitive demodulation techniques are then used to rectify the signal and determine the direction of the acceleration.\n\nInput Vector and Part Orientation\n\nThe input signal to the ADXL203 is not a standard current or voltage. Instead, the accelerometer uses the force of gravity as an input vector to determine the orientation of an object in space. Figure 2 shows the ADXL203 in five different orientations with respect to the earth’s surface and the corresponding output voltages based on the orientation of the sensor.\n\nWhen the axis of interest (the X-axis for this example) is oriented parallel to the Earth’s surface, the sensor experiences a 0 g field, which corresponds to a zero-g bias level of 2.5 V. The output voltage will change according to the sensitivity of the device (1000 mV/g). Therefore, rotating 90° clockwise (counterclockwise) will produce a +1 g field (−1 g field) and corresponding output voltage of 3.5 V (1.5 V). For various IC orientations and their associated output voltages, see Figure 2.\n\nIn order to process the accelerometer data and calculate an angle, the information must be digitized by the AD7887. It is necessary to determine the ADXL203 worst-case output voltage range and compare it to the ADC input voltage range. The AD7887 has an input voltage range of (0 V to VDD = 3.3 V). The ideal ADXL203 output voltage range is (1.5 V to 3.5 V). However, several non-idealities have been neglected in determining this range.\n\nThe first non-ideal characteristic is the zero-g bias level. This voltage is specified for 2.4 V to 2.6 V, a worst-case shift of 100 mV up or down. The second non-ideal characteristic is the sensitivity of a particular output, with a worst-case specification of 960 mV/g to 1040 mV/g. By combining both of these errors, the ADXL203 worst-case output range can be calculated:\n\nVMAX (+1 g) = (2.6 V) + (1040 mV/g)×(+1 g) = 3.64 V\n\nVMIN (−1 g) = (2.4 V) + (1040 mV/g)×(−1 g) = 1.36 V\n\nNow that the accelerometer output range has been determined, the objective is to manipulate this range (1.36 V to 3.64 V with VCM = 2.5 V) to fit the ADC input range. For dual channel operation, the AD7887 input range is 0 V to VDD (0 V to 3.3 V with VCM = 1.7 V). The quad AD8608 is used to create a 2-stage conditioning circuit as seen in Figure 1.\n\nThe first stage provides a signal gain of 1.2 and level shifts the common-mode voltage to 2 V. The second stage provides a signal gain of 1.1 (for a total signal gain of 1.32) and establishes the common-mode output voltage of 1.7 V. The output voltage range of this op amp stage lines up nicely with the ADC input voltage range, leaving approximately 200 mV headroom on the negative end and 100 mV on the positive end.\n\nSingle Axis Tilt Calculation\n\nAs an example, consider a single axis solution as indicated in Figure 3. Referring to trigonometry, the projection of the gravity vector on the X axis produces output acceleration equal to the sine of the angle between the accelerometer X axis and the horizon. The horizon is typically taken to be the plane orthogonal to the gravity vector. For an ideal value of 1 g for gravity, the output acceleration is:\n\nAX, OUT [g] = 1 g × sinθ\n\nConversion from acceleration to an inclination angle is done using the inverse sine function.\n\nθ = sin-1 (AX, OUT [g]/ g)\n\nwhere the inclination angle, θ, is in radians.\n\nIt is important to note when using a single axis solution, the sensitivity decreases as the angle between the horizon and the X axis increases. The sensitivity approaches zero as the angle approaches ±90°. This can be seen in Figure 4, where the output acceleration in g is plotted against the angle of inclination. Near ±90°, a large change in inclination angle results in a small change in output acceleration.\n\nIt is important to be cautious of signals that are out of range. It is possible for the accelerometer to output signals greater than ±1 g due to vibration, shock, or other sudden accelerations.", null, "Figure 4. Output Acceleration vs. Angle of Inclination for Single Axis Inclination Sensing\n\nSingle Axis vs. Dual Axis Considerations\n\nA simple way of addressing the decreasing sensitivity of a single axis solution as it rotates through 90° is to include a second axis perpendicular to the original. There are three major benefits to including a second axis in determining the angle of inclination.\n\nThe second major benefit of using at least two axes is that unlike the single axis solution, where tilt in any other axis will cause significant error, the use of a second axis allows an accurate value to be measured even when inclination in the third axis is present. This is because the sensitivity is proportional to the root-sum-square (rss) value of gravity on the axes of interest.", null, "Figure 6. Output Acceleration vs. Angle of Inclination for Dual-Axis Inclination Sensing\n\nThe third major benefit of using a secondary axis is the ability to distinguish between each quadrant and to measure the angles throughout an entire 360° arc. Each quadrant has a different combination of signs associated with the X and Y axis acceleration.\n\nThe inverse tangent function returns a value in Quadrant I if the operand, AX, OUT/ AY, OUT is positive; if the operand is negative, the inverse tangent function returns a value in Quadrant IV. Because the operand in Quadrant II is negative, a value of 180° should be added to the result of the calculation when the angle is in that quadrant. Because the operand in Quadrant III is positive, a value of 180° should be subtracted from the result of the calculation when the angle is in that quadrant. The correct quadrant of the calculated angle can be determined by examining the sign of the measured acceleration on each axis.\n\nDual Axis Tilt Calculation\n\nNow that a second axis is included in the system, the inclination angle calculation also requires a second look. The simple approach is to calculate the X axis as before, and to calculate the Y axis in a similar fashion, remembering to use the cosine of the angle.\n\nAX, OUT [g] = 1 g × sin θ\nAY, OUT [g] = 1 g × cos θ\n\nConvert from acceleration to an angle using the inverse sine and cosine functions.\n\nθ = sin-1 (AX, OUT [g]/ 1 g)\nθ = cos-1 (AY, OUT [g]/ 1 g)\n\nwhere the inclination angle, θ, is in radians.\n\nIt is easier, however, to apply a trigonometric identity by using the ratio of the values, which results in the following:", null, "where the inclination angle, θ, is in radians.\n\nCalibration\n\nThe most critical design aspect of the accelerometer circuit in Figure 1 is the ability to calibrate the system. Without accurate calibration, proper test procedure, and setup, the system will produce errors much larger than desired. The CN0189 Labview software includes a predefined system calibration procedure. Next, we will discuss not only how to calibrate this system, but also what contributes to the errors and why the calibration is necessary.\n\nEffects of Offset Error\n\nImagine first a dual axis solution with perfect sensitivity but with a 50 mg offset on the X axis. At 0° the X axis reads 50 mg, and the Y axis reads 1 g. The resulting calculated angle would be 2.9°, resulting in an error of 2.9°. At ±180° the X axis would report 50 mg, whereas the Y axis would report −1 g. This would result in a calculated angle and error of −2.9°.\n\nThe error between the calculated angle and the actual angle is show in Figure 7 for this example. The error due to an offset may not only be large compared to the desired accuracy of the system, but it can vary, thus making it difficult to simply calibrate out an error angle. This becomes more complicated when an offset for multiple axes is included.\n\nEffects of Sensitivity Mismatch Error\n\nThe main error component in a dual axis inclination sensing application is a difference in sensitivity between the axes of interest (in a single axis solution, any deviation between actual sensitivity and expected sensitivity results in an error). Because the ratio of the X and Y axes is used, most of the error is canceled if the sensitivities are the same.\n\nAs an example of accelerometer sensitivity mismatch, assume a dual axis solution with perfect offset trim, perfect sensitivity on the Y axis, and +5% sensitivity on the X axis. In a 1 g field, the Y axis reports 1 g and the X axis reports 1.05 g. Figure 8 shows the error in the calculated angle due to this sensitivity mismatch. Similar to offset error, the error due to accelerometer sensitivity mismatch varies over the entire range of rotation, making it difficult to compensate for the error after calculation of the inclination angle.", null, "Figure 8. Calculated Angle Error Due to Accelerometer Sensitivity Mismatch\n\nNo-Turn Calibration Technique\n\nWhen the errors due to offset and sensitivity mismatch combine, the total error can become quite large and well beyond acceptable limits in inclination sensing applications. To reduce this error, the offset and sensitivity should be calibrated, and the calibrated output acceleration used to calculate the angle of inclination. When including the effects of offset and sensitivity, the accelerometer output is as follows:\n\nAOUT[g] = AOFF + (Gain × AACTUAL)\n\nwhere:\n\nAOFF is the offset error, in g.\nGain is the gain of the accelerometer, ideally a value of 1.\nAACTUAL is the real acceleration acting on the accelerometer and the desired value, in g.\n\nA simple calibration method is to assume the gain is 1 and to measure the offset. This calibration then limits the accuracy of the system to the uncalibrated error in sensitivity. The simple calibration method can be done by placing the axis of interest into a 1 g field and measuring the output, which would be equal to the offset. That value should then be subtracted from the output of the accelerometer before processing the signal. This is often referred to as a no-turn or single point calibration because the typical orientation of a device puts the X and Y axes in a 0 g field. If a 3-axis device is used, at least one turn or a second point should be included for the Z axis.\n\nMultiple Turn Calibration Technique\n\nA more accurate calibration method is to use two points per axis of interest. When an axis is placed into a +1 g and −1 g field, the measured outputs are as follows:\n\nA+1g[g] = AOFF + (1 g × Gain)\nA−1g [g] = AOFF + (−1 g × Gain)\n\nwhere the offset, AOFF, is in 'g’.\n\nThese two points are used to determine the offset and gain as follows:\n\nAOFF [g] = 0.5 × (A+1g[g] + A-1g[g])\nGain = [0.5 × (A+1g[g] + A-1g[g])]/ 1 g\n\nwhere the +1 g and −1 g measurements, A+1g[g] and A-1g[g], are in g.\n\nThis type of calibration also helps to minimize cross-axis sensitivity effects as the orthogonal axes are in a 0 g field when making the measurements for the axis of interest. These values are used by first subtracting the offset from the accelerometer measurement and then dividing the result by the gain.\n\nAACTUAL[g] = (AOUT – AOFF) / Gain\n\nThe calculations of AOFF and Gain in the above equations assume that the acceleration values, A+1g and A−1g, are in g.\n\nIf acceleration in mg is used, the calculation of AOFF remains unchanged, but the calculation of Gain is divided by 1000 to account for the change in units.\n\nTest Results\n\nThe PCB was mounted to a board capable of spinning freely through 360°, and a set of data was taken using the calibration technique described above (finding the +1 g and −1 g values for both the X axis and Y axis to determine the offset and sensitivity of each axis). The PCB was oriented so the Y axis outputs a +1 g voltage level (~3.5 V), and the X axis outputs a 0 g voltage level (~2.5 V). This orientation, after calibration, is considered 0°.\n\nThe PCB was then rotated through ±90° in 1° increments. Figures 9 and 10 show the errors of the X and Y axis, respectively.", null, "Figure 9. Input Angle vs. Output Angle Calculated as arcsin (X)", null, "Figure 10. Input Angle vs. Output Angle Calculated as arcos (Y)", null, "Figure 11. Input Angle vs. Output Angle Calculated as arctan (X/Y)\n\nThe errors begin to increase on both axes as they approach their respective ±1 g readings. This corresponds to a board orientation of ±90° for the X axis, and 0° for the Y axis.\n\nFigure 11 shows the error according to the arctangent of the X axis and the Y axis. Notice how the error of the ratio of the two axes does not have the boundary restrictions observed in Figure 9 and Figure 10.\n\nPCB Layout Considerations\n\nIn any circuit where accuracy is crucial, it is important to consider the power supply and ground return layout on the board. The PCB should isolate the digital and analog sections as much as possible. The PCB for this system was constructed in a 4-layer stack up with large area ground plane layers and power plane polygons. See the MT-031 Tutorial for more discussion on layout and grounding and the MT-101 Tutorial for information on decoupling techniques.\n\nThe power supply to the AD7887 should be decoupled with 10 μF and 0.1 μF capacitors to properly suppress noise and reduce ripple. The capacitors should be placed as close to the device as possible with the 0.1 μF capacitor having a low ESR value. Ceramic capacitors are advised for all high frequency decoupling.\n\nPower supply lines should have as large a trace width as possible to provide low impedance paths and reduce glitch effects on the supply line. Clocks and other fast switching digital signals should be shielded from other parts of the board by digital ground.\n\nA complete design support package for this circuit can be found at www.analog.com/CN0189-DesignSupport.\n\n# Common Variations\n\nThe sensitivity of the ADXL203 and the gain of the AD7887 are both proportional to their respective supply voltages in the circuit. The entire circuit can be made ratiometric by deriving the 3.3 V VDD supply from the 5 V supply using a resistive divider followed by an AD8605 buffer as shown in Figure 12.\n\nThis configuration minimizes circuit sensitivity to supply voltage variations. The AD8505 and AD8606 are single and dual versions of the AD8608 and can be used in the circuit if desired.\n\n# Circuit Evaluation & Test\n\nThis circuit uses the EVAL-CN0189-SDPZ circuit board and the EVAL-SDP-CB1Z System Demonstration Platform (SDP) evaluation board. The two boards have 120-pin mating connectors, allowing for the quick setup and evaluation of the circuit’s performance. The EVAL-CN0189-SDPZ board contains the circuit to be evaluated, as described in this note, and the SDP evaluation board is used with the CN-0189 evaluation software to capture the data from the EVAL-CN0189-SDPZ circuit board.\n\nEquipment Needed\n\n• PC with a USB port and Windows® XP or Windows Vista® (32-bit), or Windows® 7 (32-bit)\n• EVAL-CN0189-SDPZ circuit evaluation board\n• EVAL-SDP-CB1Z SDP evaluation board\n• CN-0189 evaluation software\n• Power supply: +6 V, or +6 V “wall wart”\n\nGetting Started\n\nLoad the evaluation software by placing the CN0189 Evaluation Software disc in the CD drive of the PC. Using \"My Computer,\" locate the drive that contains the evaluation software disc and open the Readme file. Follow the instructions contained in the Readme file for installing and using the evaluation software.\n\nFunctional Block Diagram\n\nSee Figure 1 of this circuit note for the circuit block diagram, and the file “EVAL-CN0189-SDPZ-SCH-Rev0.pdf ” for the circuit schematics. This file is contained in the CN0189 Design Support Package.\n\nSetup\n\nConnect the 120-pin connector on the EVAL-CN0189-SDPZ circuit board to the connector marked “CON A” on the EVAL-SDP-CB1Z evaluation (SDP) board. Nylon hardware should be used to firmly secure the two boards, using the holes provided at the ends of the 120-pin connectors. Using an appropriate RF cable, connect the RF signal source to the EVAL-CN0189-SDPZ board via the SMA RF input connector.\n\nWith power to the supply off, connect a +6 V power supply to the pins marked “+6 V” and “GND” on the board. If available, a +6 V \"wall wart\" can be connected to the barrel connector on the board and used in place of the +6 V power supply. Connect the USB cable supplied with the SDP board to the USB port on the PC. Note: Do not connect the USB cable to the mini USB connector on the SDP board at this time.\n\nTest\n\nApply power to the +6 V supply (or “wall wart”) connected to EVAL-CN0189-SDPZ circuit board. Launch the evaluation software, and connect the USB cable from the PC to the USB mini-connector on the SDP board.\n\nOnce USB communications are established, the SDP board can be used to send, receive, and capture serial data from the EVAL-CN0189-SDPZ board.\n\nInformation and details regarding test setup, calibration, and how to use the evaluation software for data capture can be found in the CN0189 Evaluation Software Readme file.\n\nInformation regarding the SDP board can be found in the SDP User Guide." ]
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https://stat.metu.edu.tr/en/courses-0
[ "Last Updated:\n10/04/2021 - 17:08\n\nSTAT 500 Statistical Methodology in Archaeometry\n\nSubjects covering statistical methodology in collecting band analyzing data. Elementary probability distributions, hypothesis testing, analysis of variance, analysis of frequencies with emphasis on the use of computers in processing data. (Open to the students of the Archaeometry Program).\n\nSTAT 501 Statistical Theory I\n\nProbability, random variables, expectations, joint distribution functions, conditional distributions, distribution functions, moment generating functions, order statistics, censoring, limit theorems, multivariate normal distribution.\n\nSTAT 502 Statistical Theory  II\n\nLikelihood theory, sufficiency, point estimation, methods of estimation, unbiasedness, Delta method, hypothesis testing, interval estimation, asymptotic theory, Bayesian statistics, loss function, inference for bivariate distributions.\n\nSTAT 504 Nonparametric Statistical Inference and Methods\n\nUse of order statistics and other distribution-free statistics for estimation and hypothesis testing, exact non-parametric tests and measures of rank correlation. Relative efficiency, asymptotic relative efficiency and normal-score procedures. Test of goodness of fit. CCH:(1-0) 1. Prerequisite: STAT 501.\n\nSTAT 505 Sampling Theory and Methods\n\nGeneral randomization theory of simple and multistage sampling, sampling with and without replacement and with equal and unequal probabilities, ratio and regression estimates, analytical studies and multiframe problems in relation to stratification, systematic sampling, clustering and double sampling. CCH: (1-0) 1.\n\nSTAT 509 Applied Stochastic Processes\n\nMarkov chains, discrete and continuous Markov processes and associated limit theorems. Poisson and birth and death processes. Renewal processes, martingales, Brownian motion, branching processes. Weakly and strongly stationary processes, spectral analysis. Gaussian systems. CCH:(1-0)1.\n\nSTAT 518 Statistical Analysis of Designed Experiments\n\nRandomization theory of experimental design. Principles of blocking. General analysis of experimental design models. Construction and analysis of balanced and partially balanced complete and incomplete block designs. Factorial design: confounding, aliasing, fractional replication. Designs for special situations. CCH: (1-0)1.\n\nPrerequisite: STAT 501 and STAT 503.\n\nSTAT 525 Regression Theory and Methods\n\nGeneral regression models, residual analysis, selection of regression models, response surface methods, nonlinear regression models, experimental design and analysis of covariance models. Least squares, Gauss-Markov theorem. Confidence, prediction and tolerance intervals. Simultaneous inference, multiple comparison procedures. CCH: (1-0)1.\n\nSTAT 551 Probability and Statistics I\n\nProbability, combinatorics, random variables, expectations, joint distribution functions, conditional distributions, distribution functions, moment generating functions, limit theorems.\n\nSTAT 552 Probability and Statistics II\n\nOrder statistics, exponential families, sufficiency, point estimation, hypothesis testing, interval estimation, confidence intervals.\n\nSTAT 553 Actuarial Analysis and Risk Theory\n\nBasics of insurance; Basics of reinsurance; Non-life insurance mathematics; Insurance economics; Risk theory; Individual and collective risk models; Ruin theory; Credibility theory and applications.\n\nSTAT 554 Computational Statistics\n\nOverview of statistical distributions, generating random variables, exploratory data analysis, Monte Carlo (MC) method for statistical inference, data partitioning, resampling, bootstrapping, nonparametric density estimation.\n\nBivariate and multivariate smoothing, discovering structure in data, nonparametric regression, Markov Chain Monte Carlo (MCMC), statistical pattern recognition: classifiers and clustering.\n\nSTAT 556 Advanced Computing Method in Statistics\n\nThis course introduces a range of computational techniques that are important to Statistics. The topics covered include introduction to statistical computing, computer arithmetic, numerical linear algebra, regression computations, eigenproblems, numerical optimization, numerical approximations, numerical integration, expectation-maximization (EM) algorithm, basic simulation methodology, Monte Carlo (MC) integration, MC Markov Chain (MCMC) methods.\n\nSTAT 557 Statistical Modeling I\n\nIntroduction to the general theory of linear models, least squares and maximum likelihood estimation. Introduction to non-linear, log-linear and generalized linear models. Logistic and Poisson regression, ordinal and multinomial logit models. ANOVA. Causation versus association. Introduction to special Statistical Models, such as Time Series Models, Actuarial Models, Survival Models, Reliability Models.\n\nSTAT 558 Statistical Modeling II\n\nBayesian models, hierarchical modeling, nonparametric regression models, semi- parametric models, random and mixed models, response surface methods, residual analysis, correlation analysis, experimental design and analysis of covariance models.\n\nSTAT 559 Applied Multivariate Analysis\n\nCharacterizing and displaying multivariate data, multivariate distributions, tests of mean vectors and covariate matrices, discriminant analysis, classification and pattern recognition, canonical correlation, principle component analysis, factor and cluster analysis, multivariate linear, random and mixed models, multidimensional scaling.\n\nSTAT 560 Logistic Regression Analysis\n\nIntroduction to categorical response data. Fitting logistic regression models. Interpretation of coefficients. Maximum likelihood estimation. Hypothesis testing. Model building and diagnostics. Polytomous logistic regression. Interaction and confounding. Logistic regression modelling for different sampling designs: case-control and cohort studies, complex surveys. Conditional logistic regression. Exact methods for small samples. Power and sample size. Recent developments in logistic regression approach.\n\nSTAT 561 Panel Data Analysis\n\nIntroduction to longitudinal / panel data. Missing cases in panel data. Exploratory longitudinal data analysis. Marginal models, transition models, random effects models, multilevel (hierarchical) models. Estimation methods for this type of data.\n\nSTAT 562 Univariate Time Series Analysis\n\nFundamental concepts in univariate time domain analyses, properties of autocovariance and autocorrelation of time series, stationary and nonstationary models, difference equations, autoregressive integrated moving average processes, model identification, parameter estimation, model selection, time series forecasting, seasonal time series models, testing for a unit root, intervention analysis, outlier detection, handling missing observations in time series, Fourier series, spectral theory of stationary processes and the estimation of the spectrum.\n\nSTAT 563 Multivariate Time Series Analysis\n\nTransfer function models and cross-spectral analysis, time series regression and GARCH models, vector time series models, error-correction models, cointegration and causality, state space models and Kalman filter, long memory processes, nonlinear processes, temporal aggregation and disaggregation.\n\nSTAT 564 Advanced Statistical Data Analysis\n\nIntroduction to methods for analyzing experimental and observational data. Useful display of univariate and multivariate data. Exploratory data analysis. Transforming data. Detecting and handling outliers. Examining residuals. Resistant lines. Robust estimation. Approaches to handling missing data. Analysis of categorical data. Data mining.\n\nSTAT 565 Decision Theory and Bayesian Analysis\n\nIntroduction to decision making. Subjective and frequentist probability. Bayes theorem and Bayesian decision theory. Advantages of using a Bayesian approach. Likelihood principle, prior and posterior distributions, conjugate families. Inference as a statistical decision problem. Bayesian point estimation, Tests and confidence regions, model choice, invariance, equivariant estimators, hierarchical and empirical Bayes extensions, robustness and sensitivity, utility and loss, sequential experiments, Markov Chain Monte Carlo Methods, Metropolis-Hastings Algorithm, Gibbs Sampling, The E-M Algorithm.\n\nSTAT 566 Reliabiliy Theory and Methods\n\nIntroduction to reliability, order statistics, censoring and likelihood, nonparametric estimation, extreme value theory, failure time distributions, parametric likelihood concepts, simulation-based methods, testing reliability hypothesis, system reliability, failure-time regression analysis, accelerated life testing.\n\nSTAT 567 Biostatistics and Statistical Genetics\n\nIntroduction to use of statistical methodology in health related sciences. Types of health data. Odds ratio, relative risk. Prospective and retrospective study designs. Cohort, case-control, matching case-control, case-cohort, nested case-control studies. Analysis of survival data. Kaplan-Meier, life tables, Cox�s proportional hazards model. Analysis of case-control data. Unconditional, conditional, polytomous logistic regression. Introduction to genetic epidemiology. Testing Hardy-Weinberg law. Linkeage analysis. Analysis of microarray data. Association studies. Sample size and power. Recent developments in biostatistics and genetic epidemiology.\n\nSTAT 568 Statistical Consulting\n\nKey aspects of statistical consulting and data analysis activities. Formulation of statistical problems from client information. Analysis of complex data sets. Case studies. Writing and presenting reports.\n\nSTAT 542 Seminar I (0-2) NC\n\nSeminar course for M.S. students in Statistics.\n\nSTAT 543 Seminar II (0-2) NC\n\nSeminar course for M.S. students in Statistics.\n\nSTAT 544 Graduate Seminar I (0-2) NC\n\nM.S. students prepare and present a seminar in their thesis topic.\n\nSTAT 598 Term Project in Statistics (0-2) NC\n\nA project is carried out under the supervision of a faculty member in a specified area of Statistics. Students are required to write a report about their work.\n\nSTAT 599 Term Project in Interdisciplinary Statistics (0-2) NC\n\nA project is carried out under the supervision of a faculty member and an advisor from the discipline in Interdisciplinary Statistics. Students are required to write a report about their work.\n\nSTAT 601 Advanced Probability Theory I (3-0)3\n\nNotions of measure theory. General concepts and tools of probability theory. Independence; convergence; laws of large number. Random walks. Prerequisite: Consent of instructor.\n\nSTAT 602 Advanced Probability Theory II (3-0)3\n\nConcept of conditioning. From independence to dependence. Ergodic theorems. Martingales and decomposibility. Brownian motion and limit distributions. Prerequisite: Consent of instructor.\n\nSTAT 603 Advanced Theory of Statistics I (3-0)3\n\nAdvanced topics in linear and non-linear statistical estimation. Prerequisite: Consent of instructor.\n\nSTAT 604 Advanced Theory of Statistics II (3-0)3\n\nAdvanced topics in statistical hypothesis testing. Prerequisite: Consent of instructor.\n\nSTAT 605 Theory of Linear and Nonlinear Statistical Models (3-0)3\n\nGeneral linear and nonlinear models. Topics related to the statistical inference in model building. Prerequisite: Consent of instructor.\n\nSTAT 606 Theory of Experimental Designs (3-0)3\n\nBalanced and partially balanced incomplete block designs. Mixture designs. Factorial designs. Response surfaces. Optimal allocation of observations. Prerequisite: Consent of instructor.\n\nSTAT 607 Nonparametric Theory of Statistics (3-0)3\n\nRank testing and estimation procedures. Locally the most powerful rank tests. Criteria for unbiasedness. Exact and asymptotic distribution theory. Asymptotic efficiency. Rank correlation. Sequential procedures. Prerequisite: Consent of instructor.\n\nSTAT 608 Probability Models and Stochastic Processes (3-0)3\n\nDiscrete and continuous time Markov chains and Brownian motion. Gaussian processes, queues, epidemic models, branching processes, renewal processes. Prerequisite: Consent of instructor.\n\nSTAT 609 Statistical Decision Theory (3-0)3\n\nDecision theoretic approach to statistical problems. Complete class theorems. Bayes and minimax procedures. Multiple, sequential, invariant statistical decision problems. Prerequisite: Consent of instructor.\n\nSTAT 610 Sequential Analysis (3-0)3\n\nSequential probability ratio test. Approximations for stopping boundaries. Power curve and expected stopping time. Wald's lemmas. Bayes character of SPRT. Composite hypothesis. Ranking and selection CCH: (1-0)1. Prerequisite. Consent of instructor.\n\nSTAT 611 Multivariate Analysis (3-0)3\n\nAdvanced topics in multivariate statistical analysis. CCH: (1-0)1 Prerequisite: Consent of instructor.\n\nSTAT 612 Advanced Topics in Time Series Analysis (3-0)3\n\nUnivariate and multivariate time series analysis. Estimation and hypothesis testing in the time and frequency domains. CCH: (1-0)1. Prerequisite: Consent of instructor.\n\nSTAT 613 Advanced Topics in Life Testing and Reliability (3-0)3\n\nAdvanced topics in life models, reliability and hazard functions. Decision making in life testing. Design of experiments in life testing. CCH:(1-0)1. Prerequisite: Consent of instructor.\n\nSTAT 614 Interpretation of Data I (3-0)3\n\nApplication of statistical theory and procedures to various types of data. Use of computers and numerical methods are emphasized. CCH: (1-0)1. Prerequisite: Consent of instructor.\n\nSTAT 615 Interpretation of Data II (3-0)3\n\nContinuation of Stat. 614 CCH: (1-0)1. Prerequisite: Consent of instructor.\n\nSTAT 616 Applications of Statistics in Industry (3-0)3\n\nA strong background in control charts including adoptations, acceptance sampling for attributes and variables data. Acceptance plans. Statistics of combinations. CCH: (1-0)1. Prerequisite: Consent of instructor.\n\nSTAT 617 Large Sample Theory of Statistics (3-0)3\n\nLarge sample properties of tests and estimates. Problems of consistency and various forms of asymptotic efficiencies. Irregular estimation problems. Inference from stochastic processes. CCH: (1-0)1. Prerequisite: Consent of instructor.\n\nSTAT 618 Mathematical Models and Response Surface Methodology (3-0)3\n\nTwo level factorial and fractional factorial designs, blocking, polynomial models, first order and second order designs, several responses, determination and optimum conditions, design criteria involving variance and bias. CCH: (1-0)1. Prerequisite: Consent of instructor.\n\nSTAT 619 Advanced Topics in Regression and Analysis of Variance (3-0)3\n\nDevelopment of linear classification models, components of variance for balanced designs, polynomial models, harmonic regression, crossed models for combined qualitative and quantitative factors. Analysis of variance for fixed, random and mixed effects models. Randomization. Violation of assumptions. CCH: (1-0)1. Prerequisite: Consent of instructor.\n\nSTAT 620 Bayesian Inference (3-0)3\n\nSampling theory, subjective probability, likelihood principles. Bayes theorem, Bayesian analysis of normal theory, inference problems, assessment of model assumptions, robustness of inference, analysis of variance, some aspects of multivariate problems. Bayesian aspects of statistical modelling. CCH: (1-0) 1. Prerequisite: Consent of instructor.\n\nSTAT 621 Robust Statistics (3-0)3\n\nTransforming data. More refined estimators. Comparing location estimators. M and L estimators. Robust scale estimators and confidence intervals. Relevance to hypothesis testing. CCH: (1-0)1. Prerequisite: Consent of instructor.\n\nSTAT 622 Discrete Multivariate Analysis (3-0)3\n\nStructural models for counted data, maximum likelihood estimates for complete tables, formal goodness of fit; summary statistics and model selection, maximum likelihood estimates for incomplete tables, estimating the size of a closed population, models for measuring change, analysis of square tables; symmetry and marginal homogeneity, measures of association and agreement, Pseudo-Bayes estimates of cell probabilities, asymptotic methods. CCH: (1-0)1. Prerequisite: Consent of instructor.\n\nSTAT 623 Spatial Statistics (3-0)3\n\nPurely spatial processes. Spatial autocorrelation. Distribution theory for spatial statistics. Analysis for point patterns. Parametric spatial models. Estimation and testing procedures. CCH: (1-0)1. Prerequisite: Consent of instructor.\n\nSTAT 630 Advanced Topics in Statistical Inference (3-0)3\n\nSeveral advanced topics of statistical inference suited to the needs of researcher. Prerequisite: Consent of instructor.\n\nSTAT 632 Inference for Stochastic Processes (3-0)3\n\nSpecial models. Large sample theory for discrete and continuous parameter stochastic pocesses. Optimal testing. Bayesian, nonparametric and sequential inference for stochastic processes. Martingales. Stochastic differential equations. Prerequisite: Consent of instructor.\n\nSTAT 634 Theory of Stationary Random Functions (3-0)3\n\nSecond moment models of random variables and vectors. Correlation theory of random processes in the time and frequency domains. Theory of random fields in the time and frequency domains. Crossings and extremes of random functions. Applications. Prerequisite: Consent of instructor.\n\nSTAT 730 Statistics for Bioinformatics (3-0) 3\n\nDefinition of certain fundamental biological and chemical processes, principles of probability and statistics, microarray analyses, fundamental and advanced classification and clustering methods, analyses of protein sequence alignments, structure and elements of biological network, visualization tools and databases in bioinformatics.\n\nPrerequisite: Consent of instructor.\n\nSTAT 642 Seminar in Statistics I (0-2) NC\n\nSeminar course for Ph.D. students in Statistics.\n\nSTAT 643 Seminar in Statistics II (0-2) NC\n\nSeminar course for Ph.D. students in Statistics.\n\nSTAT 644 Graduate Seminar II (0-2) NC\n\nPh.D. students prepare and present a seminar in their thesis topic.\n\nSeminar course for Ph.D. students in Statistics\n\nSTAT 699 Ph.D. Thesis in Statistics (Non-credit)\n\nSTAT 800-899 Special Studies (4-2)Non-credit\n\nSTAT 900-999 Special Topics (4-0)Non-credit" ]
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https://dralb.com/2020/01/18/compositions-of-isomorphisms/
[ "# Compositions of Isomorphisms\n\nHere we will look at the composition of mappings and different properties preserved through these compositions. In particular, we will show that the composition of two isomorphisms results in an isomorphism. That is, the composition will preserve respect of operations, one-to-one and onto.\n\n## Compositions of Isomorphisms\n\nLet $$G,H$$ and $$K$$ be groups. Furthermore, let $$\\alpha:G \\to H$$ and $$\\beta:H \\to K$$ be group isomorphisms. Then $$\\beta \\circ \\alpha:G \\to K$$ is a group isomorphism.\n\n### The Theorem\n\nWith the theorem given, we want to make sure that we understand what is being stated. The first thing we want to recall is that an isomorphism is bijective homomorphism. That is, it is one-to-one and onto and respects the operation of the given group. That is, for $$\\alpha$$ we have the following,\n\n• For all $$g \\in G$$, $$\\alpha(g) \\in H$$.\n• For all $$g_{1},g_{2} \\in G$$, we have that $$\\alpha(g_{1} \\cdot_{G} g_{2})=\\alpha(g_{1}) \\cdot_{H} \\alpha(g_{2})$$.\n• For all $$h \\in H$$, there exists a $$g \\in G$$ such that $$\\alpha(g)=h$$.\n• For any $$g_{1},g_{2} \\in G$$, if $$\\alpha(g_{1})=\\alpha(g_{2})$$, then $$g_{1}=g_{2}$$.\n\nFurthermore, we would have the same properties for $$\\beta$$ if we replace the corresponding domains and codomain.\n\nThen, we must show that $$\\beta \\circ \\alpha: G \\to K$$ is an isomorphism. That is, we must show the same properties we had given above, using $$\\beta \\circ \\alpha$$ instead of $$\\alpha$$ and $$K$$ instead of $$H$$. Instead of trying to do all of these at once we will focus on one at a time.\n\n### Well-defined\n\nHere we note that we proved that the composition of any two mappings, $$\\gamma:A \\to A$$, $$\\delta:A \\to A$$, is again a well-defined mapping. Despite having different domains and codomains, the fact that the composition is well-defined follows from an extremely similar proof. We will, therefore, give the adjusted proof below.\n\n#### Proof\n\nLet $$\\alpha: G \\to H$$ and $$\\beta:H \\to K$$ be mappings and $$g \\in G$$. Then note that, since $$g \\in G$$, we have that $$\\alpha(g) \\in H$$ because $$\\alpha$$ is a well-defined mapping. Furthermore, since $$\\alpha(G) \\in H$$ and $$\\beta$$ is a well-defined mapping, we get that $$\\beta(\\alpha(g)) \\in K$$. Hence, $$\\beta \\circ \\alpha: G \\to K$$ is a well-defined mapping.\n\n#### Extra Note As an extra note, this proof goes further than needed for this problem. That is, it shows the more general case that the composition of any two mappings, $$\\alpha$$ and $$\\beta$$, if the codomain of $$\\alpha$$ is equal to the domain of $$\\beta$$, then the composition is well-defined from the domain of $$\\alpha$$ to the codomain of $$\\beta$$. ### Respects the Operation We now need to show the operation is respected for $$\\beta \\circ \\alpha$$. That is, we need to show that, if $$g_{1},g_{2} \\in G$$, then $$\\beta( \\alpha(g_{1} \\cdot_{G} g_{2}))=\\beta(\\alpha(g_{1})) \\cdot_{K} \\beta(\\alpha(g_{2}))$$. Therefore, we will start with noting \\begin{align*} \\beta(\\alpha(g_{1} \\cdot_{G} g_{2}))&=\\beta(\\alpha(g_{1})\\cdot_{H}\\alpha(g_{2})), \\end{align*} since $$\\alpha$$ is a homomorphism from $$G \\to H$$. Next, we have \\begin{align*} \\beta(\\alpha(g_{1} \\cdot_{G} g_{2}))&=\\beta(\\alpha(g_{1})\\cdot_{H}\\alpha(g_{2})) \\\\ &=\\beta(\\alpha(g_{1})) \\cdot_{K} \\beta(\\alpha(g_{2})), \\end{align*} since $$\\beta$$ is a homomorphism from $$H \\to K$$ and $$\\alpha(g_{1}),\\alpha(g_{2}) \\in H$$. We are now ready for our proof. Proof Let $$\\alpha:G \\to H$$ and $$\\beta:H \\to K$$ be group homomorphisms. Furthermore, let $$g_{1},g_{2} \\in G$$. Then \\begin{align*} \\beta(\\alpha(g_{1} \\cdot_{G} g_{2}))&=\\beta(\\alpha(g_{1})\\cdot_{H}\\alpha(g_{2})) \\\\ &=\\beta(\\alpha(g_{1})) \\cdot_{K} \\beta(\\alpha(g_{2})). \\end{align*} Hence $$\\beta \\circ \\alpha:G \\to K$$ is a group homomorphism. Extra Note As we did last, we actually proved a more general case than we need and showed that the composition of two group homomorphism is again a group homomorphism. ### One-to-one Next, we need to show that the composition of these mappings is one-to-one. That is, if we suppose that \\begin{align*} \\beta (\\alpha(g_{1})&=\\beta(\\alpha(g_{2})), \\end{align*} for some $$g_{1},g_{2} \\in G$$, then $$g_{1}=g_{2}$$. Here, we note that, because $$\\beta$$ is one-to-one and $$\\alpha(g_{1}),\\alpha(g_{2}) \\in H$$, we must have that \\begin{align*} \\alpha(g_{1})=\\alpha(g_{2}). \\end{align*} Then, since $$\\alpha$$ is one-to-one and $$g_{1}, g_{2}$$ are in $$G$$, we must have that $$g_{1}=g_{2}$$ as required. Proof Let $$\\alpha: G \\to H$$ and $$\\beta:G \\to H$$ be one-to-one mappings and $$g_{1},g_{2} \\in G$$. Furthermore, if \\)\\beta(\\alpha(g_{1}))=\\beta(\\alpha(g_{2}))\\), we have that $$\\alpha(g_{1})=\\alpha(g_{2})$$ since $$\\beta$$ is one-to-one and $$\\alpha(g_{1}),\\alpha(g_{2}) \\in H$$. Now, $$g_{1}=g_{2}$$, since $$\\alpha$$ is one-to-one and $$g_{1},g_{2} \\in G$$. Hence, $$\\beta \\circ \\alpha$$ is one-to-one. ### Onto As the last step in our proof, we need to show that $$\\beta \\circ \\alpha$$ is onto. That is, for every $$k \\in K$$, there exists a $$g \\in G$$ such that $$\\beta \\circ \\alpha(g)=k$$. Here, we can note that since $$\\beta$$ is onto, there exists an $$h \\in H$$ such that $$\\beta(h)=k$$. Furthermore, since $$\\alpha$$ is onto, there exists a $$g \\in G$$ such that $$\\alpha(g)=h$$. Then, for this $$g$$, we find that $$\\beta(\\alpha(g))=\\beta(h)=k$$. Proof Let $$\\alpha:G \\to H$$ and $$\\beta:H \\to K$$ be onto mappings of groups and $$k \\in K$$. Then, since $$\\beta$$ is onto, there exists an $$h \\in H$$ such that $$\\beta(h)=k$$. Furthermore, since $$\\alpha$$ is onto, there exists a $$g \\in G$$ such that $$\\alpha(g)=h$$. Hence, there exists a $$g \\in G$$ such that $$\\beta(\\alpha(g))=k$$, so $$\\beta \\circ \\alpha$$ is onto. ### Final Proof We now have all the components we need to show that $$\\beta \\circ \\alpha:G \\to K$$ is an isomorphism. Indeed, a formal proof would just involve placing the smaller proofs we have into one larger proof. Therefore, I will leave this to the reader. ## Conclusion Thank you for reading, and I hope this helped you while working with isomorphisms. If it did, be sure to like the post and share with someone else that may benefit from it. ### Share this: FacebookTwitterPinterestPrintEmailMoreLinkedInRedditTumblrPocketTelegramWhatsAppSkype### Like this: Like Loading... ### Related\n\nThis site uses Akismet to reduce spam. Learn how your comment data is processed." ]
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https://www.boost.org/doc/libs/1_65_1/libs/test/doc/html/boost_test/testing_tools/extended_comparison/floating_point/customizing_for_tolerance.html
[ "#", null, "Boost C++ Libraries\n\n...one of the most highly regarded and expertly designed C++ library projects in the world.\n\nThis is the documentation for an old version of Boost. Click here to view this page for the latest version.\n##### Enabling tolerance for user-defined types\n\nThe Unit Test Framework recognizes that a given type `T` is suitable for tolerance-based comparisons using the expression `boost::math::fpc::tolerance_based``<T>::value`. This meta-function already returns true for built-in floating-point types as well as any other types that match the following compile-time expression:\n\n```boost::is_floating_point<T>::value ||\n( std::numeric_limits<T>::is_specialized &&\n!std::numeric_limits<T>::is_integer &&\n!std::numeric_limits<T>::is_exact)\n```\n\nIf you require your type to also participate in tolerance-based comparisons, regardless of the above expression, you can just specialize `boost::math::fpc::tolerance_based` for your type directly, and derive it from `boost::true_type`. Your type does not even have to be a floating-point type provided that it models concept `ToleranceCompatible`.\n\n###### Example: adapting user-defined types for tolerance-based comparison\n\nCode\n\n```#define BOOST_TEST_MODULE tolerance_04\n#include <boost/test/included/unit_test.hpp>\n#include <boost/rational.hpp>\nnamespace utf = boost::unit_test;\nnamespace tt = boost::test_tools;\n\nnamespace boost { namespace math { namespace fpc {\n\ntemplate <typename I>\nstruct tolerance_based< rational<I> > : boost::true_type{};\n\n} } }\n\ntypedef boost::rational<int> ratio;\n\nBOOST_AUTO_TEST_CASE(test1, * utf::tolerance(ratio(1, 1000)))\n{\nratio x (1002, 100); // 10.02\nratio y (1001, 100); // 10.01\nratio z (1000, 100); // 10.00\n\nBOOST_TEST(x == y); // irrelevant diff by default\nBOOST_TEST(x == y, tt::tolerance(ratio(1, 2000)));\n\nBOOST_TEST(x != z); // relevant diff by default\nBOOST_TEST(x != z, tt::tolerance(ratio(2, 1000)));\n}\n```\n\nOutput\n\n```> tolerance_04\nRunning 1 test case...\ntest.cpp(23): error: in \"test1\": check x == y has failed [501/50 != 1001/100]. Relative difference exceeds tolerance [1/1001 > 1/2000]\ntest.cpp(26): error: in \"test1\": check x != z has failed [501/50 == 10/1]. Relative difference is within tolerance [1/501 < 1/500]\n\n*** 2 failures are detected in the test module \"tolerance_04\"\n```\n\n#### Concept `ToleranceCompatible`\n\n##### Refinement of\n\n`MoveConstructible`, `EqualityComparable`, `LessThanComparable`\n\n##### Notation\n\n`T`\n\nA type that is a model of `ToleranceCompatible`\n\n`x`, `y`\n\nobjects of type `T`\n\n`i`, `j`\n\nobjects of type `int`\n\n##### Valid expressions\n\nName\n\nExpression\n\nReturn type\n\nConversion from `int`\n\n```T j = i;```\n\n```x + y```\n\n`T`\n\nSubtraction\n\n```x - y```\n\n`T`\n\nNegation\n\n`-x`\n\n`T`\n\nMultiplication\n\n```x * y```\n```x * i```\n\n`T`\n\nDivision\n\n```x / y```\n```x / i```\n\n`T`\n\nMixed equality\n\n```x == i```\n```x != i```\n\n`bool`\n\nMixed ordering\n\n```x < i```\n```x > i```\n```x <= i```\n```x >= i```\n\n`bool`\n\n##### Invariants\n `T` and `int` consistency ```(x == T(i)) == (x == i)``` ```(x != T(i)) == (x != i)``` ```(x < T(i)) == (x < i)``` ```(x > T(i)) == (x > i)``` ```(x / T(i)) == (x / i)``` ```(x * T(i)) == (x * i)```" ]
[ null, "https://www.boost.org/gfx/space.png", null ]
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https://numberworld.info/201022202
[ "# Number 201022202\n\n### Properties of number 201022202\n\nCross Sum:\nFactorization:\nDivisors:\nCount of divisors:\nSum of divisors:\nPrime number?\nNo\nFibonacci number?\nNo\nBell Number?\nNo\nCatalan Number?\nNo\nBase 2 (Binary):\nBase 3 (Ternary):\nBase 4 (Quaternary):\nBase 5 (Quintal):\nBase 8 (Octal):\nbfb5afa\nBase 32:\n5vmmnq\nsin(201022202)\n0.68233394366809\ncos(201022202)\n0.73104062084015\ntan(201022202)\n0.93337350102914\nln(201022202)\n19.118925917635\nlg(201022202)\n8.3032440259459\nsqrt(201022202)\n14178.229861305\nSquare(201022202)\n\n### Number Look Up\n\nLook Up\n\n201022202 which is pronounced (two hundred one million twenty-two thousand two hundred two) is a very great figure. The cross sum of 201022202 is 11. If you factorisate 201022202 you will get these result 2 * 100511101. The figure 201022202 has 4 divisors ( 1, 2, 100511101, 201022202 ) whith a sum of 301533306. The figure 201022202 is not a prime number. The figure 201022202 is not a fibonacci number. The figure 201022202 is not a Bell Number. The figure 201022202 is not a Catalan Number. The convertion of 201022202 to base 2 (Binary) is 1011111110110101101011111010. The convertion of 201022202 to base 3 (Ternary) is 112000020222121202. The convertion of 201022202 to base 4 (Quaternary) is 23332311223322. The convertion of 201022202 to base 5 (Quintal) is 402430202302. The convertion of 201022202 to base 8 (Octal) is 1376655372. The convertion of 201022202 to base 16 (Hexadecimal) is bfb5afa. The convertion of 201022202 to base 32 is 5vmmnq. The sine of the figure 201022202 is 0.68233394366809. The cosine of the figure 201022202 is 0.73104062084015. The tangent of 201022202 is 0.93337350102914. The square root of 201022202 is 14178.229861305.\nIf you square 201022202 you will get the following result 40409925696928804. The natural logarithm of 201022202 is 19.118925917635 and the decimal logarithm is 8.3032440259459. I hope that you now know that 201022202 is unique number!" ]
[ null ]
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http://www.thermopedia.com/cn/content/1163/
[ "SUBLIMATION\n\nSublimation, or volatization, is the process of changing from a solid phase to a gaseous one, without first forming a liquid. Sublimation is one type of vaporization (see Vapor-liquid equilibrium). As with evaporation, sublimation is possible within the whole range of temperatures T and pressures p over which the solid and gaseous phases coexist. Figure 1 presents a typical phase diagram in p-T coordinates (a. water, b. carbon dioxide). It is well known that any substance can exist in one of the three states of aggregation: solid, liquid or gas. Two phase conditions can correspond to the solid state: crystal and amorphous; therefore, the notion \"phase condition\" is broader than the \"aggregate\" one. Below, however, the term \"phase transition\" implies exactly the change of the state of aggregation.", null, "Figure 1. Phase diagrams for water (a) and carbon dioxide (b).\n\nThe curves of phase equilibrium on the p-T plane intersect at the triple point, where all three states of aggregation of the substance (solid, liquid and gas) take place simultaneously. The change from a solid state to a liquid state is called \"melting\": the process of changing from a solid state to a gaseous one is called \"sublimation\" and from a liquid to a gaseous one is called \"evaporation\". The reverse process to evaporation and sublimation is called \"condensation\". The pressure", null, "at which the gaseous and condensed (liquid or solid) phases coexist is called the \"saturated vapor pressure\". For any substance relation between", null, "and T is close to exponential (the Clapeyron-CIausius Equation):", null, "where ΔQv is the heat of sublimation,", null, "vapor molecular mass, R the universal gas constant, and k is the experimentally defined constant. The heat of sublimation depends weakly on the temperature Tw.\n\nAccording to the molecular-kinetic concept, sublimation and evaporation are continuous processes of molecular emission from the interface between the gas and condensed phases, the rate of emission being governed by the thermal motion of molecules. The velocity of the reverse process (condensation) is proportional to the number of molecules per unit volume, i.e., to the partial pressure pv of the molecular species condensing on the interface. In sublimation (evaporation), a state of dynamic eqilibrium is established in a closed cavity when the condensation rate is equal to the sublimation rate. The appropriate, partial pressure is called the saturated vapor pressure, pv =", null, "(T).\n\nAccording to this model, the mass flow rate of substance during sublimation is the result for two counter processes, i.e, it is defined by the difference between the saturated vapor pressure", null, "which applies at the interface, and the partial pressure in the bulk vapor, pv the interface temperature\n\n(1)", null, "This relation is known as the Knudsen-Langmuir equation. The factor a is called the evaporation coefficient (see Accommodation Coefficient). More accurate investigations based on the methods of the molecular-kinetic theory of gases, show that in Eq. (1) the coefficient before the brackets is in the form 2a/(2 − a). This takes into account the transverse constituent of the mass velocity in the function of distribution of gas molecule velocity near the evaporation surface.\n\nWith gas flow around bodies, the process of sublimation of their surfaces is nonequilibrium. This is due to the diffusional and convective entrainment of sublimation products into the external flow (Figure 2). To predict the partial pressure pi of the ith component of the gas mixture at the surface of the body, one should consider the mass balance and allow for convective and diffusional transfer. If the mass loss rate per unit area of the body surface is Gw, and if the fraction of the ith component in the subliming material is φi, then", null, "Here, β is the mass transfer coefficient and ci is the concentration of the ith component in the boundary layer of the incoming flow. The indices w and e refer to the body surface and the boundary layer external limit.", null, "Figure 2. Sublimation from the surface of a body in crossflow.\n\nAccording to the heat transfer analogy β = (α/cp)w. To a first approximation the heat transfer coefficient (α/cp) on the sublimating smface is related to (α/cp)0 on a nonpermeable (heated) surface by the following relation:", null, "where", null, "= Gw/(α/cp)0. If we assume that in the external (oncoming) flow the products of sublimation are absent ci,e = 0, and that the sublimating body does not contain extraneous admixtures φi = 1, then we obtain the following equation for the mass loss rate:\n\n(2)", null, "This equation takes into account that the partial pressure pv is related with mass concentration cv by the relationship", null, "wherein pe, M are the pressure and molecular mass of the gas mixture.\n\nSolving Eqs. (1) and (2) simultaneously, one can reach a number of interesting conclusions. Thus, eliminating the mass loss rate Gw we obtain the relationship for estimating the degree of non-equilibrium of the sublimation process, i.e., the relation between the partial pressure pv and the saturated vapour pressure", null, ":", null, "The larger the ratio (", null, "/pv), the further the process departs from equilibrium. When the temperature and heat transfer coefficient increase, and when the pressure pe decreases, the departure from equilibrium becomes more significant. During sublimation in a vacuum, as numerous investigations of intense evaporation show, in the steady state case the maximum flow rate Gw,max =", null, "(Tw)/", null, "is not obtained, since a portion of outgoing molecules, even in the case of evaporating in a vacuum, return to the surface as a result of intermolecular collisions. It may be shown that in this case the following expression can be used (assuming a = 1):", null, "The higher the partial pressure pe of the subliming material, the closer the sublimation regime is to an equilibrium one, and closes the vapor pressure is to the saturation pressure", null, "(Tw). In this case, it can be easily shown at large values of the mass loss rate, the sublimating surface temperature asymptotically tends to a limiting value which depends only on the partial pressure in the gas flow (pe) and is defined from the equation:", null, "Introduction of the limiting temperature Tw,max is very important for estimating the possibilities of evaporative cooling (see Evaporative Cooling).\n\nUnder actual conditions, there can be a mixture of condensed media on the sublimation surface. In this case, the equation of the Knudsen-Langmuir kind is applied to each component separately. The mass loss rate is determined by summation of sublimation velocities of all components: Gw = ∑iGwi." ]
[ null, "http://www.thermopedia.com/content/5584/SUBLIMATION_FIG1.gif", null, "http://www.thermopedia.com/content/5584/eqn433.gif", null, "http://www.thermopedia.com/content/5584/eqn433.gif", null, "http://www.thermopedia.com/content/5584/eqn435.gif", null, "http://www.thermopedia.com/content/5584/eqn436.gif", null, "http://www.thermopedia.com/content/5584/eqn433.gif", null, "http://www.thermopedia.com/content/5584/eqn433.gif", null, "http://www.thermopedia.com/content/5584/eqn439.gif", null, "http://www.thermopedia.com/content/5584/eqn440.gif", null, "http://www.thermopedia.com/content/5584/SUBLIMATION_FIG2.gif", null, "http://www.thermopedia.com/content/5584/eqn441.gif", null, "http://www.thermopedia.com/content/5584/eqn442.gif", null, "http://www.thermopedia.com/content/5584/eqn443.gif", null, "http://www.thermopedia.com/content/5584/eqn444.gif", null, "http://www.thermopedia.com/content/5584/eqn433.gif", null, "http://www.thermopedia.com/content/5584/eqn446.gif", null, "http://www.thermopedia.com/content/5584/eqn433.gif", null, "http://www.thermopedia.com/content/5584/eqn433.gif", null, "http://www.thermopedia.com/content/5584/eqn449.gif", null, "http://www.thermopedia.com/content/5584/eqn450.gif", null, "http://www.thermopedia.com/content/5584/eqn433.gif", null, "http://www.thermopedia.com/content/5584/eqn452.gif", null ]
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https://byjus.com/physics/behavior-of-gases-molecules/
[ "", null, "Behavior of gas molecules and ideal gas equation\n\nWe have read about the physical properties of matter in our previous classes that included their flowability, their shape and density. As we know, the molecules of gas are loosely packed, having the least density and the greatest inter molecular distance as compared to solids and liquids.", null, "Because of this, the mutual interaction between two molecules of a gas is negligible and at low pressure and temperature, they satisfy the relation between pressure, temperature and volume as given by the following equation.\n\n$PV=KT$\n\nWhere P is the pressure of the gas, V is the volume of the gas, T is the temperature at which the system is kept, measured in Kelvin and K is the proportionality constant which depends on the nature of gases contained in the system. Further, K is proportional to the number of molecules in the sample, so we can write this mathematically as,\n\n$\\frac{PV}{KT}=constant=k_{b}$\n\nWhere N is the number of molecules in the sample and k is another constant of proportionality. The studies performed on k suggest that the value of k remains the same for all the gasses and is defined as the Boltzmann constant denoted as kb.\n\nHere we can see that,\n\nThis implies that, if P, V and T are the same, the number of molecules in the system N is the same. This proves the Avogadro’s hypothesis that for a unit volume of gasses at a fixed temperature and pressure, the number of molecules per unit volume of the gas is the same. For 22.4 liters of gas at standard temperature and pressure, this number is equal to 6.023 × 1023, which is known as the Avogadro’s number. This amount of a substance is termed as the 1 mole of that substance i.e., 6.023 × 1023 atoms of oxygen are termed as 1 mole of the oxygen atom.\n\nThis has been further proven by the kinetic theory of gasses. As per the gas equation suggested by the kinetic theory of gasses, PV=nRT , where n is the number of moles of gasses, R is the universal constant with a value equal to 8.314 Jmol-1K-1 and T is the temperature in Kelvin." ]
[ null, "https://www.facebook.com/tr", null, "https://cdn1.byjus.com/wp-content/uploads/2018/11/physics/wp-content/uploads/2016/09/10-4.png", null ]
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https://www.physicsforums.com/threads/normal-random-variables-question.406602/
[ "# Normal Random Variables Question\n\n#### Onetimeuser\n\n1. Homework Statement\n\nProblem 1 – Normal Random Variables\n\nB) Y ~ N(300, 100). Pr (300 < Y < 320) = 0.4772\n\nD) H ~ N(4000, 25). R = f(H) = 0.5H – 60. E(R) = 1940; Var(R) = 156.25\n\nI have a problem solving these problems above...I missed the class when we covered this subject and now I am lost upon solving them.. I hope somebody can help, thanks a lot\n\nRelated Calculus and Beyond Homework Help News on Phys.org\n\nHomework Helper\nFor a normally distributed random variable,\n\n$$P(a < X < b) = P(X < b) - P(X < a)$$\n\nFor any random variable $W$, if $a, b$ are real numbers,\nand\n\n$$Z = aW + b$$\n\nthen\n\n$$E(Z) = aE(W) + b, \\quad Var(Z) = a^2 Var(W)$$\n\n(as long as the mean and variance of $W$ exist)\n\n#### Onetimeuser\n\nThanks for the reply!!\n\nHowever, when I plug it in I dnt get the right answer.... did u check if the given answer is right?\n\n### Physics Forums Values\n\nWe Value Quality\n• Topics based on mainstream science\n• Proper English grammar and spelling\nWe Value Civility\n• Positive and compassionate attitudes\n• Patience while debating\nWe Value Productivity\n• Disciplined to remain on-topic\n• Recognition of own weaknesses\n• Solo and co-op problem solving" ]
[ null ]
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https://www.baeldung.com/java-and-0xff
[ "## 1. Overview\n\n0xff is a number represented in the hexadecimal numeral system (base 16). It’s composed of two F numbers in hex. As we know, F in hex is equivalent to 1111 in the binary numeral system. So, 0xff in binary is 11111111.\n\nIn this article, we’ll discover how to use the 0xff value. In addition, we’ll see how to represent it using multiple data types and how to use it with the & operator. Finally, we’ll review some of the benefits associated with using it.\n\n## 2. Representing 0xff  With Different Data Types\n\nJava allows us to define numbers interpreted as hex (base 16) by using the 0x prefix, followed by an integer literal.\n\nThe value 0xff is equivalent to 255 in unsigned decimal, -127 in signed decimal, and 11111111 in binary.\n\nSo, if we define an int variable with a value of 0xff, since Java represents integer numbers using 32 bits, the value of 0xff is 255:\n\n``````int x = 0xff;\nassertEquals(255, x);``````\n\nHowever, if we define a byte variable with the value 0xff, since Java represents a byte using 8 bits and because a byte is a signed data type, the value of 0xff is -1:\n\n``````byte y = (byte) 0xff;\nassertEquals(-1, y);``````\n\nAs we see, when we define a byte variable with the 0xff value, we need to downcast it to a byte because the range of the byte data type is from -128 to 127.\n\n## 3. Common Usage of & 0xff Operation\n\nThe & operator performs a bitwise AND operation. The output of bitwise AND is 1 if the corresponding bits of two operands is 1. On the other hand, if either bit of the operands is 0, then the result of the corresponding bit is evaluated to 0.\n\nSince 0xff has eight ones in the last 8 bits, it makes it an identity element for the bitwise AND operation. So, if we apply the x & 0xff operation, it will give us the lowest 8 bits from x. Notice that, if the number x is less than 255, it’ll still be the same. Otherwise, it’ll be the lowest 8 bits from x.\n\nIn general, the & 0xff operation provides us with a simple way to extract the lowest 8 bits from a number. We can actually use it to extract any 8 bits we need because we can shift right any of the 8 bits we want to be the lowest bits. Then, we can extract them by applying the & 0xff operation.\n\nLet’s see an example to explain some of the benefits of using & 0xff in more detail.\n\n## 4. Extracting RGBA Color Coordinates Using & 0xff\n\nLet’s assume that we have an integer number x, stored in 32 bits, that represents a color in the RGBA system, which means that it has 8 bits for each parameter (R, G, B, and A):\n\n• R = 16 (00010000 in binary)\n• G = 57  (00111001 in binary)\n• B = 168 (10101000 in binary)\n• A = 7 (00000111 in binary)\n\nSo, x in binary would be represented as 00010000 00111001 10101000 00000111 — which is the equivalent to 272214023 in decimal.\n\nNow, we have our x value in decimal, and we want to extract the value for each parameter.\n\nAs we know, the >> operation shifts bits to the right. Therefore, when we do (10000000 00000000 >> 8), it gives us 10000000. As a result, we can extract the value of each parameter:\n\n``````int rgba = 272214023;\n\nint r = rgba >> 24 & 0xff;\nassertEquals(16, r);\n\nint g = rgba >> 16 & 0xff;\nassertEquals(57, g);\n\nint b = rgba >> 8 & 0xff;\nassertEquals(168, b);\n\nint a = rgba & 0xff;\nassertEquals(7, a);``````\n\n## 5. Conclusion\n\nIn this tutorial, we’ve discussed how the & 0xff operation effectively divides a variable in a way that leaves only the value in the last 8 bits and ignores the rest of the bits. As we’ve seen, this operation is especially helpful when we shift right a variable and need to extract the shifted bits." ]
[ null ]
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http://blog.pkh.me/p/20-templating-in-c.html
[ "A small freedom area.\n\n# Templating in C\n\nSat 28 Feb 2015\n\nprog\n\nThe question of how to do templating in C is often raised, and I couldn't find a good overview of the different approaches so I'll try to make a relatively exhaustive list here.\n\nNote: the methods presented here rely on the C preprocessor, which you can use directly yourself by calling the `cpp` command, paste your code and see the output by sending an `EOS` (pressing `control + d` typically).\n\n## Simple C macro\n\nSo the most common method is to dump some code enclosed (or not) into the ```do { ... } while (0)``` form:\n\n```#define DO_RANDOM_STUFF(type) do { \\\nint i; \\\ntype *p = buf; \\\n\\\nfor (i = 0; i < len; i++) \\\np[i] = p[i] * k; \\\n} while (0)\n```\n\n... and using it directly:\n\n```int func(void *buf, int len, float k, int request)\n{\nif (request == INT8) DO_RANDOM_STUFF(int8_t);\nelse if (request == INT16) DO_RANDOM_STUFF(int16_t);\nelse if (request == INT32) DO_RANDOM_STUFF(int32_t);\n}\n```\n\n## Simple C macro, full function\n\nIt is also common to create the whole function that way:\n\n```#define DECLARE_FUNC(n) \\\nstatic void func_##n(int##n##_t *p, int len, float k) \\\n{ \\\nint i; \\\n\\\nfor (i = 0; i < len; i++) \\\np[i] = p[i] * k; \\\n}\n\nDECLARE_FUNC(8)\nDECLARE_FUNC(16)\nDECLARE_FUNC(32)\n```\n\n... which you will then use by simply calling `func_8()`, `func_16()` and `func_32()`.\n\n## Alternative function creator\n\nSo far we used the templating just to avoid typing redundancy, but it is sometimes motivated by performance. Let's observe the following pattern:\n\n```int process_image(void *img, int width, int height, const int n)\n{\nint x, y;\n\nfor (y = 0; y < height; y++) {\nfor (x = 0; x < width; x++) {\nif (n == 0) foo(img, x, y);\nelse if (n == 1) bar(img, x, y);\nelse baz(img, x, y);\n}\n}\n}\n```\n\nWe will assume here that `foo()`, `bar()` and `baz()` functions are meant to be inlined for performance reasons (be it justified or not, we assume you do not want a per-pixel function call), so you do not want to use an array of function pointers and do `func_lut[n](img, x, y)` in the inner loop. We could also consider the scenario where the functions take a completely different set of parameters.\n\nYour compiler will sometimes be smart enough to see that the `n` check in the inner loop can instead enclose the whole logic just as if you had written this:\n\n```int process_image(void *img, int width, int height, const int n)\n{\nint x, y;\n\nif (n == 0)\nfor (y = 0; y < height; y++)\nfor (x = 0; x < width; x++)\nfoo(img, x, y);\nelse if (n == 1)\nfor (y = 0; y < height; y++)\nfor (x = 0; x < width; x++)\nbar(img, x, y);\nelse\nfor (y = 0; y < height; y++)\nfor (x = 0; x < width; x++)\nbaz(img, x, y);\n}\n```\n\nUnfortunately, it is very likely that your code is much more complex (otherwise you could even just put the 2 loops into your inner functions) which will cause your compiler not to do it.\n\nOne solution for this is to create wrapper for each scenario. It would look like this:\n\n```static inline int process_image(void *img, int width, int height, const int n)\n{\nint x, y;\n\nfor (y = 0; y < height; y++) {\nfor (x = 0; x < width; x++) {\nif (n == 0) foo(img, x, y);\nelse if (n == 1) bar(img, x, y);\nelse baz(img, x, y);\n}\n}\n}\n\nint process_image_foo(void *img, int width, int height)\n{\nreturn process_image(img, width, height, 0);\n}\n\nint process_image_bar(void *img, int width, int height)\n{\nreturn process_image(img, width, height, 1);\n}\n\nint process_image_baz(void *img, int width, int height)\n{\nreturn process_image(img, width, height, 2);\n}\n```\n\nNotice how the `process_image()` function has been marked as `static inline` in order to make sure it is fully inlined in each wrapper.\n\nBy doing so, the compiler can very easily see the dead paths (because `n` is now a constant in the wrappers) and generate 3 standalone functions with zero call nor inner condition. These functions can now be put into a function lookup table mapped to `n` (`process_image_lut[n](img, width, height)`).\n\nOne great thing about this method is that you can also do type-specific code in each of the function (by interpreting `img` as `int16_t`, `float`, `uint64_t`, etc for example).\n\nAs a side effect, you will notice that the main logic is not under a huge hardly maintainable macro, and the redundancy overhead is also very small.\n\nNote: you might want to rely on `__attribute__((always_inline))` if your compiler supports it and you want the inline to be effective at every optimization level.\n\n## Mixing full functions mechanism and macros\n\nIf the redundancy overhead in the previous example is still too much for you, you can mix it with a small macro mechanism:\n\n```static inline int process_image(void *img, int width, int height, const int n)\n{\nint x, y;\n\nfor (y = 0; y < height; y++) {\nfor (x = 0; x < width; x++) {\nif (n == 0) foo(img, x, y);\nelse if (n == 1) bar(img, x, y);\nelse baz(img, x, y);\n}\n}\n}\n\n#define DECLARE_PROCESS_IMAGE_FUNC(name, n) \\\nint process_image_##name(void *img, int width, int height) \\\n{ \\\nreturn process_image(img, width, height, n); \\\n}\n\nDECLARE_PROCESS_IMAGE_FUNC(foo, 0)\nDECLARE_PROCESS_IMAGE_FUNC(bar, 1)\nDECLARE_PROCESS_IMAGE_FUNC(baz, 2)\n```\n\nThis avoids the painful redundancy and still benefits from the same advantages as previously.\n\nFFmpeg's XBR filter and paletteuse filter are a few real cases examples.\n\n## External file\n\nOne alternative to all of this is to use external files.\n\nThe logic is pretty simple; let's start with the content of `caller.c`:\n\n```#include <stdint.h>\n\n#define TEMPLATE_U16\n#include \"evil_template.c\"\n#undef TEMPLATE_U16\n\n#define TEMPLATE_U32\n#include \"evil_template.c\"\n#undef TEMPLATE_U32\n\n#define TEMPLATE_FLT\n#include \"evil_template.c\"\n#undef TEMPLATE_FLT\n\n#define TEMPLATE_DBL\n#include \"evil_template.c\"\n#undef TEMPLATE_DBL\n```\n\nThe content of `evil_template.c` could look like something like this:\n\n```#if defined(TEMPLATE_U16)\n\n# define RENAME(N) N ## _u16\n# define TYPE uint16_t\n# define SUM_TYPE uint32_t\n# define XDIV(x, n) (((x) + ((1<<(n))-1)) >> (n))\n\n#elif defined(TEMPLATE_U32)\n\n# define RENAME(N) N ## _u32\n# define TYPE uint32_t\n# define SUM_TYPE uint64_t\n# define XDIV(x, n) (((x) + ((1<<(n))-1)) >> (n))\n\n#elif defined(TEMPLATE_FLT)\n\n# define RENAME(N) N ## _flt\n# define TYPE float\n# define SUM_TYPE float\n# define XDIV(x, n) ((x) / (float)(1<<(n)))\n\n#elif defined(TEMPLATE_DBL)\n\n# define RENAME(N) N ## _dbl\n# define TYPE double\n# define SUM_TYPE double\n# define XDIV(x, n) ((x) / (double)(1<<(n)))\n\n#endif\n\nTYPE RENAME(func)(const TYPE *p, int n)\n{\nint i;\nSUM_TYPE sum = 0;\n\nfor (i = 0; i < 1<<n; i++)\nsum += p[i];\nreturn XDIV(sum, n);\n}\n\n#undef RENAME\n#undef TYPE\n#undef SUM_TYPE\n#undef XDIV\n```\n\nThis will produce the following functions:\n\n```% gcc -Wall -c caller.c -o caller.o && readelf -s caller.o | grep func\n9: 0000000000000000 110 FUNC GLOBAL DEFAULT 1 func_u16\n10: 000000000000006e 119 FUNC GLOBAL DEFAULT 1 func_u32\n11: 00000000000000e5 128 FUNC GLOBAL DEFAULT 1 func_flt\n12: 0000000000000165 131 FUNC GLOBAL DEFAULT 1 func_dbl\n```\n\nThis way of templating could be relatively handy if you have a set of functions in the template but can be considered by many as the root of all evils. Nevertheless, this is still a few order of magnitude better than C++ templating.\n\nNote: make sure your build system properly handles the dependency of `evil_template.c` from `caller.c`.\n\nResampling code in libswresample (which inspired this example) is one of the many real case example.\n\n## Conclusion\n\nMixing the full functions mechanism and macros covers most of the cases and can lead to very decent code with high level of readability and maintenance.\n\nFor more complex cases, you also can simply rely on another language, generally a scripting one (but you can use C to generate C as well of course). This can often be the most appropriate approach if the complexity is increasing and your project already has a dependency on such a language.\n\nEdit: As pointed by someone on HN, it is also worth mentioning `__Generic` in C11 that can be used for similar purpose if your project can afford the dependency on C11.\n\nindex" ]
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https://www.angelone.in/calculators/bob-fd-calculator
[ "# BOB FD Calculator\n\nTotal Investment\n\nInterest rate\n\n%\n\n1%\n\n15%\n\nSelect Duration\n\nYrs\n\n1 Yr\n\n30 Yrs\n\nThe future value of investment will be\n\nInvested Amount\n\n0\n\nEstimated Return\n\n0\n\nWhat is the BOB FD calculator?\n\nThe BOB FD calculator is a financial tool that calculates the maturity value of fixed deposit (FD) investments. Fixed deposits are one of the most popular and secure investment options, and the BOB FD calculator helps you estimate your potential returns.\n\nYou can use the BOB FD calculator to compare various investment options, such as various FD schemes with different durations and select the option with the highest returns. Customers can maximise their returns by adjusting the investment amount, tenure, or interest rate with this tool.\n\nThe BOB FD calculator is a useful tool for those who are looking to invest in fixed deposits and want to calculate the potential returns. It is simple to use and provides fast and accurate results, making it a valuable part of the FD investment process. You can find the online BOB FD calculator on Angel One and access it anytime for free.\n\nHow Does the BOB FD Calculator Work?\n\nA calculator for fixed deposits uses the compound interest formula. It uses variables such as the principal amount, tenure, and interest rate by the bank to estimate the interest and maturity proceeds. You only need to input the variables into the calculator. The online calculator quickly estimates interest and maturity proceeds.\n\nWhat is the BOB FD Calculator Formula?\n\nThe BOB FD calculator formula determines the maturity value of a fixed deposit. It considers the principal amount, interest rate, and tenure of the FD.\n\nHere’s the BOB FD formula:\n\nMaturity amount = p(1+r/n)^nt\n\nWhere,\n\np = principal amount\n\nr = rate of interest\n\nn = frequency of compounding\n\nt = tenure\n\nHow To Use the BOB FD Calculator Online?\n\nUsing the online BOB Bank FD calculator is simple and quick. Follow these steps to use the online tool:\n\n1. Enter the principal amount that you are looking to invest in the FD.\n2. Then choose the tenure or duration of the FD.\n3. Select the expected rate of interest on the fixed deposit.\n\nYou can quickly see how much the estimated interest and maturity amount is.\n\nLet’s understand the above steps with an example. Suppose you want to invest Rs. 80,000 in a fixed deposit for 5 years at 7.5% p.a. Using the online FD calculator, calculate the interest and maturity proceeds as follows:\n\n1. Enter Rs. 80,000 in the principal field of the tool\n2. Key in the interest rate as 7.5% p.a.\n3. Enter 5 years as the tenure\n4. Click on ‘Calculate’\n\nThe FD calculator will quickly calculate the interest and maturity amount. In this case, the interest earned will be Rs. 35,996. The maturity amount will be Rs. 115,996.\n\nBenefits of Using BOB FD Calculator\n• Informed investments - One of the main benefits of using the BOB FD calculator is that it allows investors to make informed investment decisions. Investors can calculate their expected returns at maturity by entering the principal amount, tenure, and interest rates. They can also experiment with various combinations of principal amount, tenure, and interest rate to find the best investment option.\n• Quick and error-free - Another benefit of using the BOB FD calculator is that it helps you save time and effort. Investors can use the calculator instead of manually calculating the returns on their fixed deposits to get accurate results quickly. This not only saves time but also ensures that the calculations are error-free.\n• Convenient - The BOB FD calculator is also a simple tool that can be found online. Investors can use it from the convenience of their own homes, eliminating the need to visit a bank branch or speak with a customer service representative\n• ## FAQs\n\n### What is the online BOB FD calculator?", null, "The BOB FD calculator is an online tool that calculates the interest earned and maturity amount for Bank of Baroda fixed deposits. Users can get instant results by entering the deposit amount, tenure, and interest rate.\n\n### How to use the BOB FD calculator online?", null, "Navigate to the FD calculator on the Angel One website and enter the deposit amount, tenure, and interest rate. The calculator will show the maturity amount as well as the interest earned on your investment.\n\n### Is the BOB calculator free to use?", null, "Yes, the BOB calculator is completely free to use. There are no hidden fees or charges associated with using it. Simply visit the Angel One website and start calculating.\n\n### How is fixed deposit interest calculated?", null, "The fixed deposit interest amount can be calculated by using an online FD calculator to find out the maturity amount for a particular principal. The difference between the principal and the maturity amount is the interest amount.\n\n### What is the interest rate on BOB FD?", null, "Bank of Baroda (BoB) offers fixed deposit (FD) at attractive interest rates. The rates depend on the tenure and principal amount. The higher the variables, the higher the returns.\n\nEnjoy Zero Brokerage on\nEquity Delivery\n\n##### Enjoy Zero Brokerage on Equity Delivery", null, "" ]
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http://benjaminfarnell.com/gfnmj/c36213-is-the-square-root-of-50-a-whole-number
[ "For example: 16 divided by 4 is 4. The only square root of zero is zero. The square root of 49 is 7 and the square room of 64 is 8, so all of the number between 49 and 64 have a square root that is a decimal between 7 and 8. Some examples are: 2, 9, 20417 and 32. For example 25 is a perfect square since $$\\pm \\sqrt{25}= \\pm 5$$ If the radicand is not a perfect square i.e. 1. Irrational Square Root. (a) Does the square root of the number 40 lie between the 6 and 7? I got some good news, and some bad news. STEP 6: Subtract Again. Negative numbers don't have real square roots since a square is either positive or 0. In mathematics, a square root of a number x is a number y such that y 2 = x; in other words, a number y whose square (the result of multiplying the number by itself, or y ⋅ y) is x. The good news is that the square root of a whole number is rational precisely when it is an integer. Can the square root of a real number be negative? i.e., the square root of 4 is √4=2 which when multiplied by itself gives the original definite number 4. I know 78 lies between 64 and 81 (both perfect squares). Square Root Table: In mathematics, the square of a number refers to the value which we get after multiplying the same number by itself (Y × Y = X). That is, it is the number under the root symbol. 1, … For example, 2 is the square root of 4, because 2x2=4. Some common roots include the square root, where n = 2, and the cubed root, where n = 3. For example, use the square root calculator below to find the square root of 5. A number is a perfect square when you can directly extract its square root in whole numbers and without approximation. You want to know between which two numbers the square root of 78 lies. it's square root is 5, since 5 x 5 is equal to 25. Write 5 next to 4 in the top right corner. Square root of a number is a value, which on multiplied by itself gives the original number. Thus, a = 8 and b = 9. Go here for the next problem on our list. To find the square root of a whole number, you could also divide the whole number by numbers until you get an answer that is the same as the number you used to divide the whole number. Simplifying the square root of a whole number. I.e., a and b such that a < sqrt(78) < b. How can I check if the number produced after finding the square root of a number is a whole number? What would happen if $\\sqrt{}$ instead meant the negative square root? A perfect square root is where the square root of a number equals another whole number. Is it wrong to say that $100$ is solution of $\\sqrt x+10=0$? 2 & 3 & 4 because the square root of 2 is 4, the square root of 3 is 9, and the square root of 4 is 16. Only numbers bigger than or equal to zero have real square roots. The reason the square root of 0.2003 is greater than 0.2003 is because when you take the square root of the dividend (√2003), the decrease of the dividend is smaller than the decrease of the divisor when you take the square root of the divisor (√10000). Unlike the examples above, not every square root of a number ends up being a nice and neat whole number. A number bigger than zero has two square roots: one is positive (bigger than zero) and the other is negative (smaller than zero). Edited February 2, 2007 by BALA We will find a whole number bigger than the square root of 102 and a whole number smaller than the square root of 102. Example #2: Estimate the square root of 102. It is the second digit in the root. In equation format: n √ a = b b n = a. Estimating a Root. For numbers that are perfect squares, you can find whole numbers as answers. Natural (Counting) Numbers: Whole Numbers: Natural Numbers and . Well, we can just consider a^2 < 78 < b^2. (b) Identify all whole numbers whose square root lie between 6 and 7. This would not be the case if the whole number in front of the decimal point wasn't 0. Use division to find the square root. Example: 25 is a perfect square. Example: 1 - 1 4 - 2 9 - 3 16 - 4... up to 961 - 31...Which is the last square root before 1000. It is an irrational number, but you can simplify it or find rational approximations for it.. First note that #50 = 2 xx 5 xx 5# contains a square factor #5^2#.We can use this to simplify the square root: #sqrt(50) = sqrt(5^2*2) = sqrt(5^2)*sqrt(2) = 5 sqrt(2)# Many square roots of numbers turn out to be irrational roots, that is irrational numbers. Give reasons in support of your answer. Find the square root of 56.25, without using long division method? These are equivalent statements, but both of them are irrational. To square root is to find the two identical factors of a number. Determine the Type of Number square root of 49. The square root of 17 is between 4 and 5. Which is both a real number and an integer? For the numbers above, the square root was equal to an integer. When you subject a certain number to a radical sign, you can extract its square. Note that the term \"radicand\" refers to the number for which the root is to be determined. Sometimes, you may get a real number when finding the square root. 5.858585858 63.4 square root 21 square root 36 2. Examples. They are all integers, which means they do not have decimal points or fractions. For example, 4 has two square roots: 2 and -2. Simplifying the square root of a whole number. 55. Here, the square root of X (√X) refers to Y.Every non-negative number such as 1,2,3,4,5,…, etc., can have a non-negative square root such as √4=2,√9=3,√16=4, etc.The square root lists can be written in a table. 2. We could estimate the square of 17 to be 4.1 for example. * List of Perfect Squares At first glance, this would appear to be so, because the poster's example finds the square root of the two digit whole number 20 instead of the article's example of 645. This is the same thing as the square root of 1/100. However, I actually worked out the article's example (square root of 645) using both methods and found that the Babylonian Method required 9 \"cycles of divide and average\" to arrive at the answer. There are six common sets of numbers. “Is the square root of 20 an integer and/or a rational, irrational, whole, or natural number?” I have the unsettling feeling that I’m being asked to do someone’s homework for them. Subtract the product we calculated (which is 425) from the current number on the left (also 425). Add the approximate root with the original number divided by the approximate root and divide by 2. x_i := (x_i + n / x_i) / 2; Continue step 2 until the difference in the approximate root along … Which of these numbers can be classified as both real and irrational? Recall also that natural numbers are positive real whole numbers (i.e. Check out the Wikipedia pages on square root and methods of computing square roots. You take an irrational number, divide it by 10, you still have an irrational number. Prime Factors can help determine if a number will have a square root that is rational or irrational. I.e., sqrt(78) lies between 8 and 9. Irrational Numbers: Non Terminating or Non Repeating Decimals. The perfect squares are the squares of the whole numbers: 1, 4, 9, 16, 25, 36, 49, 64, 81, 100 … The square root of a number, n, written . Square root confusion? This is an online chart which consists of a list of square root values for numbers 1 to 100. Here are the square roots of all the perfect squares from 1 to 100. How to find the square root of a number. You could also use you calculate and hit the square root key and then type 51 and see what decimal you get: about 7.1, which is between 7 and 8. Square Root of 50 in Decimal form rounded to nearest 5 decimals: 7.07107 Exponent Form Square Root of 50 written with Exponent instead of Radical: 50 ½ = 5 x 2 ½ Simplify Square Root of 51 The answer to Simplify Square Root of 50 is not the only problem we solved. Suppose, x is the square root of y, then it is represented as x=√y or we can express the same equation as x 2 = y. Here,’√’is the radical symbol used to represent the root of numbers. √ All whole numbers lie to the right of the zero on the number line. 85 times 5 results in 425, which is exactly what we need. The square root of #50# is not a whole number, or even a rational number. And 4 divided by 2 is 2, and so on. is the number that gives n when multiplied by itself. The foot of a ladder is placed 6 feet away from a wall. Integers: Rational Numbers: Integers, Fractions, and Terminating or Repeating Decimals. But I’ll bite. it is not always possible to get the square root as an integer. The number 8 and the number 5 give us 85. Calculating square roots and n th roots is fairly intensive. For example, because 10 x 10 = 100. ... Square root of whole number number of solutions. However, for numbers that aren't perfect squares, you'll have to use a method that involves estimation (or you can use a table of square and square roots). If the square root of an integer is another integer then the square is called a perfect square. 1. For example, the square root of 144 is 12. In mathematics, the general root, or the n th root of a number a is another number b that when multiplied by itself n times, equals a. To take the square root of a number, press [2ND] (the secondary function key) and then [√ ] (the radical symbol key which is used to take the square root of a number) and then the number that you want to find the square root of and then the [ENTER] key.Example: To find the square root of 2, push: [2ND] [√ ] 2 [ENTER] This will give you the answer of: 1.414213562 if done correctly. 12 is a whole number thus 144 is a perfect square root. A square root is number which produces a definite number when multiplied by itself. We have the square root of 0.01. And they are all whole numbers. Which is equal to the square root of 1 over the square root of 100, which is equal to 1/10, or 0.1. Take a reasonable guess (approximate root) for the square root. 23 1 over 4 square root 27 3.402538 3. The square root of 5 is 25 and the square root of 6 is 36 so they would not belong in the particular catagory you're asking for. I need to Output the results of square root just for the whole numbers. Let's do J." ]
[ null ]
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https://answers.everydaycalculation.com/divide-fractions/1-75-divided-by-6-9
[ "# Answers\n\nSolutions by everydaycalculation.com\n\n## Divide 1/75 with 6/9\n\n1/75 ÷ 6/9 is 1/50.\n\n#### Steps for dividing fractions\n\n1. Find the reciprocal of the divisor\nReciprocal of 6/9: 9/6\n2. Now, multiply it with the dividend\nSo, 1/75 ÷ 6/9 = 1/75 × 9/6\n3. = 1 × 9/75 × 6 = 9/450\n4. After reducing the fraction, the answer is 1/50\n\nMathStep (Works offline)", null, "Download our mobile app and learn to work with fractions in your own time:\nAndroid and iPhone/ iPad\n\n#### Divide Fractions Calculator\n\n÷\n\n© everydaycalculation.com" ]
[ null, "https://answers.everydaycalculation.com/mathstep-app-icon.png", null ]
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https://catchsomeair.us/and-relationship/what-is-the-relationship-among-electric-force-charge-and-distance.php
[ "# What is the relationship among electric force charge and distance\n\n### Comparing electric force and gravitational force (practice) | Khan Academy", null, "Practice: Relationship between electric force, charge, and distance analyzing the electric and gravitational interaction between a proton and an alpha particle. only really need to do the Qs, to see the charge (Q1 multiplied by Q2) E=k( Q1Q1/r^2), describe the relationship between electric force, charge, and distance. Forces between two electrically-charged objects can be extremely large. Metals are good conductors of electric charge, while plastics, wood, and rubber are not. When the ground connection is removed, the conductor will have a charge opposite in sign to that r is the distance between the charges.\n\nThe magnitude of the force and the distance between the two balloons is said to be inversely related.\n\n### Electric charge and Coulomb's law\n\nCoulomb's Law Equation The quantitative expression for the effect of these three variables on electric force is known as Coulomb's law.\n\nCoulomb's law states that the electrical force between two charged objects is directly proportional to the product of the quantity of charge on the objects and inversely proportional to the square of the separation distance between the two objects. In equation form, Coulomb's law can be stated as where Q1 represents the quantity of charge on object 1 in CoulombsQ2 represents the quantity of charge on object 2 in Coulombsand d represents the distance of separation between the two objects in meters.\n\nThe symbol k is a proportionality constant known as the Coulomb's law constant. The value of this constant is dependent upon the medium that the charged objects are immersed in.", null, "In the case of air, the value is approximately 9. If the charged objects are present in water, the value of k can be reduced by as much as a factor of It is worthwhile to point out that the units on k are such that when substituted into the equation the units on charge Coulombs and the units on distance meters will be canceled, leaving a Newton as the unit of force. The Coulomb's law equation provides an accurate description of the force between two objects whenever the objects act as point charges.\n\nA charged conducting sphere interacts with other charged objects as though all of its charge were located at its center. While the charge is uniformly spread across the surface of the sphere, the center of charge can be considered to be the center of the sphere. The sphere acts as a point charge with its excess charge located at its center. Since Coulomb's law applies to point charges, the distance d in the equation is the distance between the centers of charge for both objects not the distance between their nearest surfaces.\n\nThe symbols Q1 and Q2 in the Coulomb's law equation represent the quantities of charge on the two interacting objects. The sign on the charge is simply representative of whether the object has an excess of electrons a negatively charged object or a shortage of electrons a positively charged object.\n\nWhile the practice is not recommended, there is certainly no harm in doing so.", null, "This is consistent with the concept that oppositely charged objects have an attractive interaction and like charged objects have a repulsive interaction. If the charge of one of the objects is doubled, and the distance separating the objects is doubled, then what is the new force? The electrostatic force is directly related to the product of the charges and inversely related to the square of the separation distance. Doubling one of the charges would serve to double the force. Doubling the distance would serve to reduce the force by a factor of four.\n\nThe combined affect of these two variations would be to decrease the force by a factor of two - changing it from 0. If the charge of both of the objects is doubled and the distance separating the objects is doubled, then what is the new force? Doubling both of the charges would serve to quadruple the force. The combined affect of these two variations would be to not change the force at all; it remains as 0.\n\nIf the charge of one of the objects is increased by a factor of four, and the distance separating the objects is doubled, then what is the new force? Quadrupling one of the charges would serve to quadruple the force. The combined affect of these two variations would be to not alter the force at all; it would remain as 0.\n\nThere are three ways that objects can be given a net charge. Charging by friction - this is useful for charging insulators. If you rub one material with another say, a plastic ruler with a piece of paper towelelectrons have a tendency to be transferred from one material to the other.\n\nFor example, rubbing glass with silk or saran wrap generally leaves the glass with a positive charge; rubbing PVC rod with fur generally gives the rod a negative charge.\n\n### Relationship between electric force, charge, and distance (practice) | Khan Academy\n\nCharging by conduction - useful for charging metals and other conductors. If a charged object touches a conductor, some charge will be transferred between the object and the conductor, charging the conductor with the same sign as the charge on the object.\n\nCharging by induction - also useful for charging metals and other conductors. Again, a charged object is used, but this time it is only brought close to the conductor, and does not touch it. If the conductor is connected to ground ground is basically anything neutral that can give up electrons to, or take electrons from, an objectelectrons will either flow on to it or away from it.", null, "When the ground connection is removedthe conductor will have a charge opposite in sign to that of the charged object. An example of induction using a negatively charged object and an initially-uncharged conductor for example, a metal ball on a plastic handle. Electrons on the conductor will be repelled from the area nearest the charged object.\n\nThe electrons on the conductor want to get as far away from the negatively-charged object as possible, so some of them flow to ground. This leaves the conductor with a deficit of electrons.\n\n## Inverse Square Law\n\nThe conductor is now positively charged. A practical application involving the transfer of charge is in how laser printers and photocopiers work. This is a good web page that gives a nice description of how a photocopier works: University of Delaware Why is static electricity more apparent in winter?\n\nYou notice static electricity much more in winter with clothes in a dryer, or taking a sweater off, or getting a shock when you touch something after walking on carpet than in summer because the air is much drier in winter than summer. Dry air is a relatively good electrical insulator, so if something is charged the charge tends to stay. In more humid conditions, such as you find on a typical summer day, water molecules, which are polarized, can quickly remove charge from a charged object.\n\nTry this at home See if you can charge something at home using friction. I got good results by rubbing a Bic pen with a piece of paper towel. To test the charge, you can use a narrow stream of water from a faucet; if the object attracts the stream when it's brought close, you know it's charged.\n\n## CHAPTER 22\n\nAll you need to do is to find something to rub - try anything made out of hard plastic or rubber. You also need to find something to rub the object with - potential candidates are things like paper towel, wool, silk, and saran wrap or other plastic.\n\nCoulomb's law The force exerted by one charge q on another charge Q is given by Coulomb's law: Remember that force is a vector, so when more than one charge exerts a force on another charge, the net force on that charge is the vector sum of the individual forces." ]
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https://en.m.wikisource.org/wiki/1911_Encyclop%C3%A6dia_Britannica/Mathematics
[ "# 1911 Encyclopædia Britannica/Mathematics\n\nMATHEMATICS (Gr. μαθηματική, sc. τέχνη or ἐπιστήμη; from μάθημα, “learning” or “science”), the general term for the various applications of mathematical thought, the traditional field of which is number and quantity. It has been usual to define mathematics as “the science of discrete and continuous magnitude.” Even Leibnitz, who initiated a more modern point of view, follows the tradition in thus confining the scope of mathematics properly so called, while apparently conceiving it as a department of a yet wider science of reasoning. A short consideration of some leading topics of the science will exemplify both the plausibility and inadequacy of the above definition. Arithmetic, algebra, and the infinitesimal calculus, are sciences directly concerned with integral numbers, rational (or fractional) numbers, and real numbers generally, which include incommensurable numbers. It would seem that “the general theory of discrete and continuous quantity” is the exact description of the topics of these sciences. Furthermore, can we not complete the circle of the mathematical sciences by adding geometry? Now geometry deals with points, lines, planes and cubic contents. Of these all except points are quantities: lines involve lengths, planes involve areas, and cubic contents involve volumes. Also, as the Cartesian geometry shows, all the relations between points are expressible in terms of geometric quantities. Accordingly, at first sight it seems reasonable to define geometry in some such way as “the science of dimensional quantity.” Thus every subdivision of mathematical science would appear to deal with quantity, and the definition of mathematics as “the science of quantity” would appear to be justified. We have now to consider the reasons for rejecting this definition as inadequate.\n\nTypes of Critical Questions.—What are numbers? We can talk of five apples and ten pears. But what are “five” and “ten” apart from the apples and pears? Also in addition to the cardinal numbers there are the ordinal numbers: the fifth apple and the tenth pear claim thought. What is the relation of “the fifth” and “the tenth” to “five” and “ten”? “The first rose of summer” and “the last rose of summer” are parallel phrases, yet one explicitly introduces an ordinal number and the other does not. Again, “half a foot” and “half a pound” are easily defined. But in what sense is there “a half,” which is the same for “half a foot” as “half a pound”? Furthermore, incommensurable numbers are defined as the limits arrived at as the result of certain procedures with rational numbers. But how do we know that there is anything to reach? We must know that √2 exists before we can prove that any procedure will reach it. An expedition to the North Pole has nothing to reach unless the earth rotates.\n\nAlso in geometry, what is a point? The straightness of a straight line and the planeness of a plane require consideration. Furthermore, “congruence” is a difficulty. For when a triangle “moves,” the points do not move with it. So what is it that keeps unaltered in the moving triangle? Thus the whole method of measurement in geometry as described in the elementary textbooks and the older treatises is obscure to the last degree. Lastly, what are “dimensions”? All these topics require thorough discussion before we can rest content with the definition of mathematics as the general science of magnitude; and by the time they are discussed the definition has evaporated. An outline of the modern answers to questions such as the above will now be given. A critical defence of them would require a volume.\n\nCardinal Numbers.—A one-one relation between the members of two classes α and β is any method of correlating all the members of α to all the members of β, so that any member of α has one and only one correlate in β, and any member of β has one and only one correlate in α. Two classes between which a one-one relation exists have the same cardinal number and are called cardinally similar; and the cardinal number of the class α is a certain class whose members are themselves classes—namely, it is the class composed of all those classes for which a one-one correlation with α exists. Thus the cardinal number of α is itself a class, and furthermore α is a member of it. For a one-one relation can be established between the members of α and α by the simple process of correlating each member of α with itself. Thus the cardinal number one is the class of unit classes, the cardinal number two is the class of doublets, and so on. Also a unit class is any class with the property that it possesses a member x such that, if y is any member of the class, then x and y are identical. A doublet is any class which possesses a member x such that the modified class formed by all the other members except x is a unit class. And so on for all the finite cardinals, which are thus defined successively. The cardinal number zero is the class of classes with no members; but there is only one such class, namely—the null class. Thus this cardinal number has only one member. The operations of addition and multiplication of two given cardinal numbers can be defined by taking two classes α and β, satisfying the conditions (1) that their cardinal numbers are respectively the given numbers, and (2) that they contain no member in common, and then by defining by reference to α and β two other suitable classes whose cardinal numbers are defined to be respectively the required sum and product of the cardinal numbers in question. We need not here consider the details of this process.\n\nWith these definitions it is now possible to prove the following six premisses applying to finite cardinal numbers, from which Peano has shown that all arithmetic can be deduced:—\n\ni. Cardinal numbers form a class.\nii. Zero is a cardinal number.\niii. If a is a cardinal number, a+1 is a cardinal number.\niv. If s is any class and zero is a member of it, also if when x is a cardinal number and a member of s, also x+1 is a member of s, then the whole class of cardinal numbers is contained in s.\nv. If a and b are cardinal numbers, and a+1 =b+1, then a = b.\nvi. If a is a cardinal number, then a+1 ≠ 0.\n\nIt may be noticed that (iv) is the familiar principle of mathematical induction. Peano in an historical note refers its first explicit employment, although without a general enunciation, to Maurolycus in his work, Arithmeticorum libri duo (Venice, 1575).\n\nBut now the difficulty of confining mathematics to being the science of number and quantity is immediately apparent. For there is no self-contained science of cardinal numbers. The proof of the six premisses requires an elaborate investigation into the general properties of classes and relations which can be deduced by the strictest reasoning from our ultimate logical principles. Also it is purely arbitrary to erect the consequences of these six principles into a separate science. They are excellent principles of the highest value, but they are in no sense the necessary premisses which must be proved before any other propositions of cardinal numbers can be established. On the contrary, the premisses of arithmetic can be put in other forms, and, furthermore, an indefinite number of propositions of arithmetic can be proved directly from logical principles without mentioning them. Thus, while arithmetic may be defined as that branch of deductive reasoning concerning classes and relations which is concerned with the establishment of propositions concerning cardinal numbers, it must be added that the introduction of cardinal numbers makes no great break in this general science. It is no more than an interesting subdivision in a general theory.\n\nOrdinal Numbers.—We must first understand what is meant by “order,” that is, by “serial arrangement.” An order of a set of things is to be sought in that relation holding between members of the set which constitutes that order. The set viewed as a class has many orders. Thus the telegraph posts along a certain road have a space-order very obvious to our senses; but they have also a time-order according to dates of erection, perhaps more important to the postal authorities who replace them after fixed intervals. A set of cardinal numbers have an order of magnitude, often called the order of the set because of its insistent obviousness to us; but, if they are the numbers drawn in a lottery, their time-order of occurrence in that drawing also ranges them in an order of some importance. Thus the order is defined by the “serial” relation. A relation (R) is serial when (i) it implies diversity, so that, if x has the relation R to y, x is diverse from y; (2) it is transitive, so that if x has the relation R to y, and у to z, then x has the relation R to z; (3) it has the property of connexity, so that if x and у are things to which any things bear the relation R, or which bear the relation R to any things, then either x is identical with y, or x has the relation R to y, or у has the relation R to x. These conditions are necessary and sufficient to secure that our ordinary ideas of “preceding” and “succeeding” hold in respect to the relation R. The “field” of the relation R is the class of things ranged in order by it. Two relations R and R′ are said to be ordinally similar, if a one-one relation holds between the members of the two fields of R and R', such that if x and у are any two members of the field of R, such that x has the relation R to y, and if x′ and y′ are the correlates in the field of R' of x and y, then in all such cases x′ has the relation R' to y′, and conversely, interchanging the dashes on the letters, i.e. R and R', x and x′, &c. It is evident that the ordinal similarity of two relations implies the cardinal similarity of their fields, but not conversely. Also, two relations need not be serial in order to be ordinally similar; but if one is serial, so is the other. The relation-number of a relation is the class whose members are all those relations which are ordinally similar to it. This class will include the original relation itself. The relation-number of a relation should be compared with the cardinal number of a class. When a relation is serial its relation-number is often called its serial type. The addition and multiplication of two relation-numbers is defined by taking two relations R and S, such that (1) their fields have no terms in common; (2) their relation-numbers are the two relation- numbers in question, and then by defining by reference to R and S two other suitable relations whose relation-numbers are defined to be respectively the sum and product of the relation-numbers in question. We need not consider the details of this process. Now if n be any finite cardinal number, it can be proved that the class of those serial relations, which have a field whose cardinal number is n, is a relation-number. This relation-number is the ordinal number corresponding to n; let it be symbolized by . Thus, corresponding to the cardinal numbers 2, 3, 4 … there are the ordinal numbers 2., 3., 4. … The definition of the ordinal number 1 requires some little ingenuity owing to the fact that no serial relation can have a field whose cardinal number is 1; but we must omit here the explanation of the process. The ordinal number 0· is the class whose sole member is the null relation—that is, the relation which never holds between any pair of entities. The definitions of the finite ordinals can be expressed without use of the corresponding cardinals, so there is no essential priority of cardinals to ordinals. Here also it can be seen that the science of the finite ordinals is a particular subdivision of the general theory of classes and relations. Thus the illusory nature of the traditional definition of mathematics is again illustrated.\n\nCantor's Infinite Numbers.—Owing to the correspondence between the finite cardinals and the finite ordinals, the propositions of cardinal arithmetic and ordinal arithmetic correspond point by point. But the definition of the cardinal number of a class applies when the class is not finite, and it can be proved that there are different infinite cardinal numbers, and that there is a least infinite cardinal, now usually denoted by א0, where א is the Hebrew letter aleph. Similarly, a class of serial relations, called well-ordered serial relations, can be defined, such that their corresponding relation-numbers include the ordinary finite ordinals, but also include relation-numbers which have many properties like those of the finite ordinals, though the fields of the relations belonging to them are not finite. These relation-numbers are the infinite ordinal numbers. The arithmetic of the infinite cardinals does not correspond to that of the infinite ordinals. The theory of these extensions of the ideas of number is dealt with in the article Number. It will suffice to mention here that Peano's fourth premiss of arithmetic does not hold for infinite cardinals or for infinite ordinals. Contrasting the above definitions of number, cardinal and ordinals, with the alternative theory that number is an ultimate idea incapable of definition, we notice that our procedure exacts a greater attention, combined with a smaller credulity; for every idea, assumed as ultimate, demands a separate act of faith.\n\nThe Data of Analysis.—Rational numbers and real numbers in general can now be defined according to the same general method. If m and n are finite cardinal numbers, the rational number m/n is the relation which any finite cardinal number x bears to any finite cardinal number y when n×x = m×y. Thus the rational number one, which we will denote by 1r, is not the cardinal number 1; for 1r, is the relation 1/1 as defined above, and is thus a relation holding between certain pairs of cardinals. Similarly, the other rational integers must be distinguished from the corresponding cardinals. The arithmetic of rational numbers is now established by means of appropriate definitions, which indicate the entities meant by the operations of addition and multiplication. But the desire to obtain general enunciations of theorems without exceptional cases has Ted mathematicians to employ entities of ever-ascending types of elaboration. These entities are not created by mathematicians, they are employed by them, and their definitions should point out the construction of the new entities in terms of those already on hand. The real numbers, which include irrational numbers, have now to be defined. Consider the serial arrangement of the rationals in their order of magnitude. A real number is a class (α, say) of rational numbers which satisfies the condition that it is the same as the class of those rationals each of which precedes at least one member of α. Thus, consider the class of rationals less than 2r; any member of this class precedes some other members of the class—thus 1/2 precedes 4/3, 3/2 and so on; also the class of predecessors of predecessors of 2r, is itself the class of predecessors of 2r. Accordingly this class is a real number; it will be called the real number 2R. Note that the class of rationals less than or equal to 2r is not a real number. For 2r is not a predecessor of some member of the class. In the above example 2R is an integral real number, which is distinct from a rational integer, and from a cardinal number. Similarly, any rational real number is distinct from the corresponding rational number. But now the irrational real numbers have all made their appearance. For example, the class of rationals whose squares are less than 2r satisfies the definition of a real number; it is the real number √2. The arithmetic of real numbers follows from appropriate definitions of the operations of addition and multiplication. Except for the immediate purposes of an explanation, such as the above, it is unnecessary for mathematicians to have separate symbols, such as 2, 2r, and 2R, or 2/3 and (2/3)R. Real numbers with signs (+ or −) are now defined. If a is a real number, +a is defined to be the relation which any real number of the form x+a bears to the real number x, and −a is the relation which any real number x bears to the real number x+a. The addition and multiplication of these “signed” real numbers is suitably defined, and it is proved that the usual arithmetic of such numbers follows. Finally, we reach a complex number of the nth order. Such a number is a “one-many” relation which relates n signed real numbers (or n algebraic complex numbers when they are already defined by this procedure) to the n cardinal numbers 1, 2 … n respectively. If such a complex number is written (as usual) in the form x1e1+x2e2+ … +xnen, then this particular complex number relates x1 to 1, x2 to 2, … x1 to n. Also the “unit” e1 (or es) considered as a number of the system is merely a shortened form for the complex number (+1)e1+0e2+. . . +0en. This last number exemplifies the fact that one signed real number, such as 0, may be correlated to many of the n cardinals, such аз 2 … n in the example, but that each cardinal is only correlated with one signed number. Hence the relation has been called above “one-many.” The sum of two complex numbers x1e1+x2e2+ … +xnen and y1e1+y2e2+ … +ynen is always defined to be the complex number (x1+y2)e1+ (x2+y2)e2+ … + (xn+yn)en. But an indefinite number of definitions of the product of two complex numbers yield interesting results. Each definition gives rise to a corresponding algebra of higher complex numbers. We will confine ourselves here to algebraic complex numbers— that is, to complex numbers of the second order taken in connexion with that definition of multiplication which leads to ordinary algebra. The product of two complex numbers of the second order—namely, x1e1+x2e2 and y1e1+y2e2, is in this case defined to mean the complex (x1y1x2y2)e1+ (x1y2+x2y1)e2. Thus e1×e1 = e1, e2×e2 = −e1, e1×e2 = e2×e1 = e2. With this definition it is usual to omit the first symbol e1, and to write i or √−1 instead of e2. Accordingly, the typical form for such a complex number is x+yi, and then with this notation the above-mentioned definition of multiplication is invariably adopted. The importance of this algebra arises from the fact that in terms of such complex numbers with this definition of multiplication the utmost generality of expression, to the exclusion of exceptional cases, can be obtained for theorems which occur in analogous forms, but complicated with exceptional cases, in the algebras of real numbers and of signed real numbers. This is exactly the same reason as that which has led mathematicians to work with signed real numbers in preference to real numbers, and with real numbers in preference to rational numbers. The evolution of mathematical thought in the invention of the data of analysis has thus been completely traced in outline.\n\nDefinition of Mathematics.—It has now become apparent that the traditional field of mathematics in the province of discrete and continuous number can only be separated from the general abstract theory of classes and relations by a wavering and indeterminate line. Of course a discussion as to the mere application of a word easily degenerates into the most fruitless logomachy. It is open to any one to use any word in any sense. But on the assumption that “mathematics” is to denote a science well marked out by its subject matter and its methods from other topics of thought, and that at least it is to include all topics habitually assigned to it, there is now no option but to employ “mathematics” in the general sense of the “science concerned with the logical deduction of consequences from the general premisses of all reasoning.”\n\nGeometry.—The typical mathematical proposition is: “If x, y, z … satisfy such and such conditions, then such and such other conditions hold with respect to them.” By taking fixed conditions for the hypothesis of such a proposition a definite department of mathematics is marked out. For example, geometry is such a department. The “axioms” of geometry are the fixed conditions which occur in the hypotheses of the geometrical propositions. The special nature of the “axioms” which constitute geometry is considered in the article Geometry (Axioms). It is sufficient to observe here that they are concerned with special types of classes of classes and of classes of relations, and that the connexion of geometry with number and magnitude is in no way an essential part of the foundation of the science. In fact, the whole theory of measurement in geometry arises at a comparatively late stage as the result of a variety of complicated considerations.\n\nClasses and Relations.—The foregoing account of the nature of mathematics necessitates a strict deduction of the general properties of classes and relations from the ultimate logical premisses. In the course of this process, undertaken for the first time with the rigour of mathematicians, some contradictions have become apparent. That first discovered is known as Burali-Forti's contradiction, and consists in the proof that there both is and is not a greatest infinite ordinal number. But these contradictions do not depend upon any theory of number, for Russell 's contradiction does not involve number in any form. This contradiction arises from considering the class possessing as members all classes which are not members of themselves. Call this class w; then to say that x is a w is equivalent to saying that x is not an x. Accordingly, to say that w is a w is equivalent to saying that w is not a w. An analogous contradiction can be found for relations. It follows that a careful scrutiny of the very idea of classes and relations is required. Note that classes are here required in extension, so that the class of human beings and the class of rational featherless bipeds are identical; similarly for relations, which are to be determined by the entities related. Now a class in respect to its components is many. In what sense then can it be one? This problem of “the one and the many” has been discussed continuously by the philosophers. All the contradictions can be avoided, and yet the use of classes and relations can be preserved as required by mathematics, and indeed by common sense, by a theory which denies to a class—or relation— existence or being in any sense in which the entities composing it— or related by it—exist. Thus, to say that a pen is an entity and the class of pens is an entity is merely a play upon the word “entity”; the second sense of “entity” (if any) is indeed derived from the first, but has a more complex signification. Consider an incomplete proposition, incomplete in the sense that some entity which ought to be involved in it is represented by an undetermined x, which may stand for any entity. Call it a prepositional function; and, if φх be a prepositional function, the undetermined variable x is the argument. Two prepositional functions φх and ψх are “extensionally identical” if any determination of x in φх which converts φх into a true proposition also converts ψх into a true proposition, and conversely for ψ and φ. Now consider a prepositional function Fχ in which the variable argument χ is itself a prepositional function. If Fχ is true when, and only when, χ is determined to be either φ or some other prepositional function extensionally equivalent to φ, then the proposition Fφ is of the form which is ordinarily recognized as being about the class determined by φх taken in extension—that is, the class of entities for which φх is a true proposition when x is determined to be any one of them. A similar theory holds for relations which arise from the consideration of prepositional functions with two or more variable arguments. It is then possible to define by a parallel elaboration what is meant by classes of classes, classes of relations, relations between classes, and so on. Accordingly, the number of a class of relations can be defined, or of a class of classes, and so on. This theory is in effect a theory of the use of classes and relations, and does not decide the philosophic question as to the sense (if any) in which a class in extension is one entity. It does indeed deny that it is an entity in the sense in which one of its members is an entity. Accordingly, it is a fallacy for any determination of x to consider “x is an x” or “x is not an x” as having the meaning of propositions. Note that for any determination of x, “x is an x” and “x is not an x,” are neither of them fallacies but are both meaningless, according to this theory. Thus Russell's contradiction vanishes, and an examination of the other contradictions shows that they vanish also.\n\nApplied Mathematics.—The selection of the topics of mathematical inquiry among the infinite variety open to it has been guided by the useful applications, and indeed the abstract theory has only recently been disentangled from the empirical elements connected with these applications. For example, the application of the theory of cardinal numbers to classes of physical entities involves in practice some process of counting. It is only recently that the succession of processes which is involved in any act of counting has been seen to be irrelevant to the idea of number. Indeed, it is only by experience that we can know that any definite process of counting will give the true cardinal number of some class of entities. It is perfectly possible to imagine a universe in which any act of counting by a being in it annihilated some members of the class counted during the time and only during the time of its continuance. A legend of the Council of Nicea illustrates this point: “When the Bishops took their places on their thrones, they were 318; when they rose up to be called over, it appeared that they were 319; so that they never could make the number come right, and whenever they approached the last of the series, he immediately turned into the likeness of his next neighbour.” Whatever be the historical worth of this story, it may safely be said that it cannot be disproved by deductive reasoning from the premisses of abstract logic. The most we can do is to assert that a universe in which such things are liable to happen on a large scale is unfitted for the practical application of the theory of cardinal numbers. The application of the theory of real numbers to physical quantities involves analogous considerations. In the first place, some physical process of addition is presupposed, involving some inductively inferred law of permanence during that process. Thus in the theory of masses we must know that two pounds of lead when put together will counterbalance in the scales two pounds of sugar, or a pound of lead and a pound of sugar. Furthermore, the sort of continuity of the series (in order of magnitude) of rational numbers is known to be different from that of the series of real numbers. Indeed, mathematicians now reserve “continuity” as the term for the latter kind of continuity; the mere property of having an infinite number of terms between any two terms is called “compactness.” The compactness of the series of rational numbers is consistent with quasi-gaps in it—that is, with the possible absence of limits to classes in it. Thus the class of rational numbers whose squares are less than 2 has no upper limit among the rational numbers. But among the real numbers all classes have limits. Now, owing to the necessary inexactness of measurement, it is impossible to discriminate directly whether any kind of continuous physical quantity possesses the compactness of the series of rationals or the continuity of the series of real numbers. In calculations the latter hypothesis is made because of its mathematical simplicity. But, the assumption has certainly no a priori grounds in its favour, and it is not very easy to see how to base it upon experience. For example, if it should turn out that the mass of a body is to be estimated by counting the number of corpuscles (whatever they may be) which go to form it, then a body with an irrational measure of mass is intrinsically impossible. Similarly, the continuity of space apparently rests upon sheer assumption unsupported by any a priori or experimental grounds. Thus the current applications of mathematics to the analysis of phenomena can be justified by no a priori necessity.\n\nIn one sense there is no science of applied mathematics. When once the fixed conditions which any hypothetical group of entities are to satisfy have been precisely formulated, the deduction of the further propositions, which also will hold respecting them, can proceed in complete independence of the question as to whether or no any such group of entities can be found in the world of phenomena. Thus rational mechanics, based on the Newtonian Laws, viewed as mathematics is independent of its supposed application, and hydrodynamics remains a coherent and respected science though it is extremely improbable that any perfect fluid exists in the physical world. But this unbendingly logical point of view cannot be the last word upon the matter. For no one can doubt the essential difference between characteristic treatises upon “pure” and “applied” mathematics. The difference is a difference in method. In pure mathematics the hypotheses which a set of entities are to satisfy are given, and a group of interesting deductions are sought. In “applied mathematics” the “deductions” are given in the shape of the experimental evidence of natural science, and the hypotheses from which the “deductions” can be deduced are sought. Accordingly, every treatise on applied mathematics, properly so-called, is directed to the criticism of the “laws” from which the reasoning starts, or to a suggestion of results which experiment may hope to find. Thus if it calculates the result of some experiment, it is not the experimentalist’s well-attested results which are on their trial, but the basis of the calculation. Newton's Hypotheses non fingo was a proud boast, but it rests upon an entire misconception of the capacities of the mind of man in dealing with external nature.\n\nSynopsis of Existing Developments of Pure Mathematics.—A complete classification of mathematical sciences, as they at present exist, is to be found in the International Catalogue of Scientific Literature promoted by the Royal Society. The classification in question was drawn up by an international committee of eminent mathematicians, and thus has the highest authority. It would be unfair to criticize it from an exacting philosophical point of view. The practical object of the enterprise required that the proportionate quantity of yearly output in the various branches, and that the liability of various topics as a matter of fact to occur in connexion with each other, should modify the classification.\n\nSynopsis of Existing Developments of Applied Mathematics.— Section В of the International Catalogue deals with mechanics. The heading “Measurement of Dynamical Quantities” includes the topics units, measurements, and the constant of gravitation. The topics of the other headings do not require express mention. These headings are: “Geometry and Kinematics of Particles and Solid Bodies”; “Principles of Rational Mechanics”; “Statics of Particles, Rigid Bodies, &c.”; “Kinetics of Particles, Rigid Bodies, &c.”; “General Analytical Mechanics”; “Statics and Dynamics of Fluids”; “Hydraulics and Fluid Resistances”; “Elasticity.” For the subjects of this general heading see the articles Mechanics; Dynamics, Analytical; Gyroscope and Gyrostat; Harmonic Analysis; Wave; Hydromechanics; Elasticity; Motion, Law of; Energy; Energetics; Astronomy (Celestial Mechanics); Tide. Mechanics (including dynamical astronomy) is that subject among those traditionally classed as “applied” which has been most completely transfused by mathematics—that is to say, which is studied with the deductive spirit of the pure mathematician, and not with the covert inductive intention overlaid with the superficial forms of deduction, characteristic of the applied mathematician.\n\nEvery branch of physics gives rise to an application of mathematics. A prophecy may be hazarded that in the future these applications will unify themselves into a mathematical theory of a hypothetical substructure of the universe, uniform under all the diverse phenomena. This reflection is suggested by the following articles: Aether; Molecule; Capillary Action; Diffusion; Radiation, Theory of; and others.\n\nThe applications of mathematics to statistics (see Statistics and Probability) should not be lost sight of; the leading fields for these applications are insurance, sociology, variation in zoology and economics.\n\nThe History of Mathematics.—The history of mathematics is in the main the history of its various branches. A short account of the history of each branch will be found in connexion with the article which deals with it. Viewing the subject as a whole, and apart from remote developments which have not in fact seriously influenced the great structure of the mathematics of the European races, it may be said to have had its origin with the Greeks, working on pre-existing fragmentary lines of thought derived from the Egyptians and Phoenicians. The Greeks created the sciences of geometry and of number as applied to the measurement of continuous quantities. The great abstract ideas (considered directly and not merely in tacit use) which have dominated the science were due to them—namely, ratio, irrationality, continuity, the point, the straight line, the plane. This period lasted from the time of Thales, с. 600 B.C., to the capture of Alexandria by the Mahommedans, A.D. 641. The medieval Arabians invented our system of numeration and developed algebra. The next period of advance stretches from the Renaissance to Newton and Leibnitz at the end of the i7th century. During this period logarithms were invented, trigonometry and algebra developed, analytical geometry invented, dynamics put upon a sound basis, and the period closed with the magnificent invention of (or at least the perfecting of) the differential calculus by Newton and Leibnitz and the discovery of gravitation. The i8th century witnessed a rapid development of analysis, and the period culminated with the genius of Lagrange and Laplace. This period may be conceived as continuing throughout the first quarter of the 19th century. It was remarkable both for the brilliance of its achievements and for the large number of French mathematicians of the first rank who flourished during it. The next period was inaugurated in analysis by K. F. Gauss, N. H. Abel and A. L. Cauchy. Between them the general theory of the complex variable, and of the various “infinite” processes of mathematical analysis, was established, while other mathematicians, such as Poncelet, Steiner, Lobatschewsky and von Staudt, were founding modern geometry, and Gauss inaugurated the differential geometry of surfaces. The applied mathematical sciences of light, electricity and electromagnetism, and of heat, were now largely developed. This school of mathematical thought lasted beyond the middle of the century, after which a change and further development can be traced. In the next and last period the progress of pure mathematics has been dominated by the critical spirit introduced by the German mathematicians under the guidance of Weierstrass, though foreshadowed by earlier analysts, such as Abel. Also such ideas as those of invariants, groups and of form, have modified the entire science. But the progress in all directions has been too rapid to admit of any one adequate characterization. During the same period a brilliant group of mathematical physicists, notably Lord Kelvin (W. Thomson), H. V. Helmholtz, J. C. Maxwell, H. Hertz, have transformed applied mathematics by systematically basing their deductions upon the Law of the conservation of energy, and the hypothesis of an ether pervading space.\n\nBibliography.—References to the works containing expositions of the various branches of mathematics are given in the appropriate articles. It must suffice here to refer to sources in which the subject is considered as one whole. Most philosophers refer in their works to mathematics more or less cursorily, either in the treatment of the ideas of number and magnitude, or in their consideration of the alleged a priori and necessary truths. A bibliography of such references would be in effect a bibliography of metaphysics, or rather of epistemology. The founder of the modern point of view, explained in this article, was Leibnitz, who, however, was so far in advance of contemporary thought that his ideas remained neglected and undeveloped until recently; cf. Opuscules et fragments inédits de Leibnitz. Extraits des manuscrits de la bibliothèque royale de Hanovre, by Louis Couturat (Paris, 1903), especially pp. 356-399, “Generales inquisitiones de analysi notionum et veritatum” (written in 1686); also cf. La Logique de Leibnitz, already referred to. For the modern authors who have rediscovered and improved upon the position of Leibnitz, cf. Grundgesetze der Arithmetik, begriffsschrifttich abgeleitet von Dr G. Frege, a. o. Professor an der Univ. Jena (Bd. i., 1893; Bd. ii., 1903, Jena); also cf. Frege's earlier works, Begriffsschrift, eine der arithmetischen nachgebildete Formelsprache des reinen Denkens (Halle, 1879), and Die Grundlagen der Anthmetik (Breslau, 1884); also cf. Bertrand Russell, The Principles of Mathematics (Cambridge, 1903), and his article on “Mathematical Logic” in Amer. Quart. Journ. of Math. (vol. xxx., 1908). Also the following works are of importance, though not all expressly expounding the Leibnitzian point of view: cf. G. Cantor, “Grundlagen einer allgemeinen Mannigfaltigkeitslehre,” Math. Annal., voi. xxi. (1883) and subsequent articles in vols. xlvi. and xlix.; also R. Dedekind, Stetigkeit und irrationales Zahlen (1st ed., 1872), and Was sind und was sollen die Zahlen? (1st ed., 1887), both tracts translated into English under the title Essays on the Theory of Numbers (Chicago, 1901). These works of G. Cantor and Dedekind were of the greatest importance in the progress of the subject. Also cf. G. Peano (with various collaborators of the Italian school), Formulaire de mathématiques (Turin, various editions, 1894–1908; the earlier editions are the more interesting philosophically); Felix Klein, Lectures on Mathematics (New York, 1894); W. K. Clifford, The Common Sense of the exact Sciences (London, 1885); H. Poincaré, La Science et l’hypothèse (Paris, ist ed., 1902), English translation under the title, Science and Hypothesis (London, 1905); L. Couturat, Les Principes des mathématiques (Paris, 1905) ; E. Mach, Die Mechanik in ihrer Entwicklung (Prague, 1883), English translation under the title, The Science of Mechanics (London, 1893); K. Pearson, The Grammar of Science (London, 1st ed., 1892; and ed., 1900, enlarged); A. Cayley, Presidential Address (Brit. Assoc., 1883); B. Russell and A. N. Whitehead, Principia Mathematica (Cambridge, 1911). For the history of mathematics the one modern and complete source of information is M. Cantor’s Vorlesungen über Geschichte der Mathematik (Leipzig, 1st Bd., 1880; 2nd Bd., 1892; 3rd Bd., 1898; 4th Bd., 1908; 1st Bd., von den ältesten Zeiten bis zum Jahre 1200, n. Chr.; 2nd Bd., von 1200–1668; 3rd Bd., von 1668–1758; 4th Bd., von 1795 bis 1799); W. W. R. Ball, A Short History of Mathematics (London 1st ed., 1888, three subsequent editions, enlarged and revised, and translations into French and Italian).  (A. N. W.)\n\n1. Cf. La Logique de Leibnitz, ch. vii., by L. Couturat (Paris, 1901).\n2. Cf. The Principles of Mathematics, by Bertrand Russell (Cambridge, 1903).\n3. Cf. Formulaire mathématique (Turin, ed. of 1903); earlier formulations of the bases of arithmetic are given by him in the editions of 1898 and of 1901. The variations are only trivial.\n4. Cf. Russell, loc. cit., pp. 199–256.\n5. The first unqualified explicit statement of part of this definition seems to be by B. Peirce, “Mathematics is the science which draws necessary conclusions” (Linear Associative Algebra, § i. (1870), republished in the Amer. Journ. of Math., vol. iy. (1881) ). But it will be noticed that the second half of the definition in the text—“from the general premisses of all reasoning”—is left unexpressed. The full expression of the idea and its development into a philosophy of mathematics is due to Russell, loc. cit.\n6. “Una questione sui numeri transfiniti,” Rend, del circolo mat. di Palermo, vol. xi. (1897); and Russell, loc. cit., ch. xxxviii.\n7. Cf. Russell, loc. cit., ch. x.\n8. Cf. Pragmatism: a New Name for some Old Ways of Thinking (1907).\n9. Due to Bertrand Russell, cf. “Mathematical Logic as based on the Theory of Types,” Amer. Journ. of Math. vol. xxx. (1908). It is more fully explained by him, with later simplifications, in Principia mathematica (Cambridge).\n10. Cf. Stanley’s Eastern Church, Lecture v.\n11. Cf. A Short History of Mathematics, by W. W. R. Ball." ]
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https://zbmath.org/?q=an:06803441
[ "## Sums of asymptotically midpoint uniformly convex spaces.(English)Zbl 1388.46018\n\nSummary: We study the property of asymptotic midpoint uniform convexity for infinite direct sums of Banach spaces, where the norm of the sum is defined by a Banach space $$E$$ with a 1-unconditional basis. We show that a sum $$(\\sum_{n=1}^\\infty X_n)_E$$ is asymptotically midpoint uniformly convex (AMUC) if and only if the spaces $$X_n$$ are uniformly AMUC and $$E$$ is uniformly monotone. We also show that $$L_p(X)$$ is AMUC if and only if $$X$$ is uniformly convex.\n\n### MSC:\n\n 46B20 Geometry and structure of normed linear spaces\n\n### Keywords:\n\nuniform convexity; asymptotic geometry; asymptotic moduli; AMUC\nFull Text:" ]
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https://unoamor.ch/3-EWIG%20UNO%20BIBLIO/MAIL-WIR-BITTEN-DRUCKEN.htm
[ ".\nU      N      O\nS  P  I  R  I  T\n\nF I R S T\nU  N  O    A  M  O  R    U  N  O\nF  L  O  R      A  Z  U  L\n\n.\nU   N   O\n\nG E B O R E N E    M E N S C H E N    L E B E    W E S E N\nA N T O N I O    G U T E R R E S   /   T A T I A N A    V A L O V A Y A\n.\nA L S    M E N S C H E N    B I T T E N    W I R\n\nA N T O N I O    G U T E R R E S    UND    T A T I A N A    V A L O V A Y A\nD I E    M Ö G L I C H K E I T    I N    B E T R A C H T    Z U    Z I E H E N\nB E I    D E R\nU   N   O   G\nE I N    G E M E I N S A M    O R T          E W I G    U N O    B I B L I O\nZ U    B E G R Ü N D E N\n.\nE I N    G E M E I N S A M    O R T\nI N D E M    W I R    U N S    A L S\nU N O    S P I R I T    I S T E N\nU N O    A M O R    I S T E N\n\nE I N    W E L T    I C H    G E I S T    B Ü R G E R\n\nU N O    U N I T E D    P E O P L E\nF Ö R D E R E R\nD U R C H    U N S E R E\nU N O  %  U N O      E I N   G A B E\nF Ü R    D I E    U N    R E F O R M    >>>\nU P    U N O    U N\nE I N    G E B E N    K Ö N N E N\n.\nW I R    B I T T E N\nA N T O N I O    G U T E R R E S    UND    T A T I A N A    V A L O V A Y A\nG E M\nE I N S A M    M I T    S I M O N E    S O M M A R U G A\n.\nD  E  N      S  A  M  E  N\nB L A U E    B L U M E    H O F F N U N G    F L O R    A Z U L\nU  N  O      A  M  O  R      U  N  O\nE  W  I  G      U  N  O      B  I  B  L  I  O\n.\nI  D  E  E  L      E   I   N      G  E  B  E  T  T  E  T\nI N    D E R\nL  I  C  H  T      E  R  D  -  S  A  M  E  N      W  A  S  S  E  R      L  U  F  T\nH  Ü  T  T  E\nA U F    D E M\nM  O  N  T  E    U   N   O    V  E  R  I  T  À\nE  I  N      Z U      S  Ä  E  N\n.\nU N D    D A S    O F F I C E\nE W I G    U N O    B I B L I O\nF  L  O  R      A  Z  U  L\nB E I    D E R    U N O G\nE  I  N      Z U      P  F  L  A  N  Z  E  N\nU N D    U N K O N V E N T I O N E L L    Z U    R E A L I S I E R E N\n.\nW I R    B I T T E N\nA N T O N I O    G U T E R R E S    UND    T A T I A N A    V A L O V A Y A\nU N   +   U N O G    V E R T R E T E R\n.\nU N S    Z U    B E N A C H R I C H T I G E N  -  W A N N    W I R    U N S    A L S\nU N O    A M O R    I S T E N\nF Ü R    D I E\nU  P      U   N   O      U N\nZ I E L E\nM  E  N  S  C  H  L  I  C  H  K  E  I  T              M  E  N  S  C  H  E  N  R  E  C  H  T\nU  N  O      P  L  A  N  E  T      S  C  H  U  T  Z\nA  B  C  P      W  A  F  F  E  N      A  B  S  C  H  A  F  F  E  N\nE  I  N      G  E  B  E  N      D  Ü  R  F  E  N\n.\nI N    D A N K B A R K E I T    G R Ü S S T    S I E    H E R Z L I C H\nE W I G    U N O    A M O R    I S T\n.\nAUS FÜLLEN     AUS DRUCKEN     UNTER SCHREIBEN     EIN SCHREIBEN     EIN SENDEN\nOFFICE UNO BIBLIO   UNOG   PALAIS DES NATIONS   CH - 1211 - GENF\n.\n\n . ORGANISATION  NAME : . ADRESSE : . MAIL ADRESSE : . UNTERSCHRIFT : ." ]
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http://isabelle.in.tum.de/repos/isabelle/file/f7750d816c21/src/HOL/HOL.thy
[ "src/HOL/HOL.thy\n author wenzelm Thu Mar 11 13:20:35 1999 +0100 (1999-03-11) changeset 6349 f7750d816c21 parent 6340 7d5cbd5819a0 child 6795 35f214e73668 permissions -rw-r--r--\nremoved foo_build_completed -- now handled by session management (via usedir);\n``` 1 (* Title: HOL/HOL.thy\n```\n``` 2 ID: \\$Id\\$\n```\n``` 3 Author: Tobias Nipkow\n```\n``` 4 Copyright 1993 University of Cambridge\n```\n``` 5\n```\n``` 6 Higher-Order Logic.\n```\n``` 7 *)\n```\n``` 8\n```\n``` 9 HOL = CPure +\n```\n``` 10\n```\n``` 11\n```\n``` 12 (** Core syntax **)\n```\n``` 13\n```\n``` 14 global\n```\n``` 15\n```\n``` 16 classes\n```\n``` 17 term < logic\n```\n``` 18\n```\n``` 19 default\n```\n``` 20 term\n```\n``` 21\n```\n``` 22 types\n```\n``` 23 bool\n```\n``` 24\n```\n``` 25 arities\n```\n``` 26 fun :: (term, term) term\n```\n``` 27 bool :: term\n```\n``` 28\n```\n``` 29\n```\n``` 30 consts\n```\n``` 31\n```\n``` 32 (* Constants *)\n```\n``` 33\n```\n``` 34 Trueprop :: bool => prop (\"(_)\" 5)\n```\n``` 35 Not :: bool => bool (\"~ _\" 40)\n```\n``` 36 True, False :: bool\n```\n``` 37 If :: [bool, 'a, 'a] => 'a (\"(if (_)/ then (_)/ else (_))\" 10)\n```\n``` 38 arbitrary :: 'a\n```\n``` 39\n```\n``` 40 (* Binders *)\n```\n``` 41\n```\n``` 42 Eps :: ('a => bool) => 'a\n```\n``` 43 All :: ('a => bool) => bool (binder \"! \" 10)\n```\n``` 44 Ex :: ('a => bool) => bool (binder \"? \" 10)\n```\n``` 45 Ex1 :: ('a => bool) => bool (binder \"?! \" 10)\n```\n``` 46 Let :: ['a, 'a => 'b] => 'b\n```\n``` 47\n```\n``` 48 (* Infixes *)\n```\n``` 49\n```\n``` 50 \"=\" :: ['a, 'a] => bool (infixl 50)\n```\n``` 51 \"&\" :: [bool, bool] => bool (infixr 35)\n```\n``` 52 \"|\" :: [bool, bool] => bool (infixr 30)\n```\n``` 53 \"-->\" :: [bool, bool] => bool (infixr 25)\n```\n``` 54\n```\n``` 55\n```\n``` 56 (* Overloaded Constants *)\n```\n``` 57\n```\n``` 58 axclass\n```\n``` 59 plus < term\n```\n``` 60\n```\n``` 61 axclass\n```\n``` 62 minus < term\n```\n``` 63\n```\n``` 64 axclass\n```\n``` 65 times < term\n```\n``` 66\n```\n``` 67 axclass\n```\n``` 68 power < term\n```\n``` 69\n```\n``` 70 consts\n```\n``` 71 \"+\" :: ['a::plus, 'a] => 'a (infixl 65)\n```\n``` 72 \"-\" :: ['a::minus, 'a] => 'a (infixl 65)\n```\n``` 73 uminus :: ['a::minus] => 'a (\"- _\" 100)\n```\n``` 74 \"*\" :: ['a::times, 'a] => 'a (infixl 70)\n```\n``` 75 (*See Nat.thy for \"^\"*)\n```\n``` 76\n```\n``` 77\n```\n``` 78 (** Additional concrete syntax **)\n```\n``` 79\n```\n``` 80 nonterminals\n```\n``` 81 letbinds letbind\n```\n``` 82 case_syn cases_syn\n```\n``` 83\n```\n``` 84 syntax\n```\n``` 85\n```\n``` 86 \"~=\" :: ['a, 'a] => bool (infixl 50)\n```\n``` 87\n```\n``` 88 \"@Eps\" :: [pttrn, bool] => 'a (\"(3@ _./ _)\" [0, 10] 10)\n```\n``` 89\n```\n``` 90 (* Alternative Quantifiers *)\n```\n``` 91\n```\n``` 92 \"*All\" :: [idts, bool] => bool (\"(3ALL _./ _)\" [0, 10] 10)\n```\n``` 93 \"*Ex\" :: [idts, bool] => bool (\"(3EX _./ _)\" [0, 10] 10)\n```\n``` 94 \"*Ex1\" :: [idts, bool] => bool (\"(3EX! _./ _)\" [0, 10] 10)\n```\n``` 95\n```\n``` 96 (* Let expressions *)\n```\n``` 97\n```\n``` 98 \"_bind\" :: [pttrn, 'a] => letbind (\"(2_ =/ _)\" 10)\n```\n``` 99 \"\" :: letbind => letbinds (\"_\")\n```\n``` 100 \"_binds\" :: [letbind, letbinds] => letbinds (\"_;/ _\")\n```\n``` 101 \"_Let\" :: [letbinds, 'a] => 'a (\"(let (_)/ in (_))\" 10)\n```\n``` 102\n```\n``` 103 (* Case expressions *)\n```\n``` 104\n```\n``` 105 \"@case\" :: ['a, cases_syn] => 'b (\"(case _ of/ _)\" 10)\n```\n``` 106 \"@case1\" :: ['a, 'b] => case_syn (\"(2_ =>/ _)\" 10)\n```\n``` 107 \"\" :: case_syn => cases_syn (\"_\")\n```\n``` 108 \"@case2\" :: [case_syn, cases_syn] => cases_syn (\"_/ | _\")\n```\n``` 109\n```\n``` 110 translations\n```\n``` 111 \"x ~= y\" == \"~ (x = y)\"\n```\n``` 112 \"@ x. b\" == \"Eps (%x. b)\"\n```\n``` 113 \"ALL xs. P\" => \"! xs. P\"\n```\n``` 114 \"EX xs. P\" => \"? xs. P\"\n```\n``` 115 \"EX! xs. P\" => \"?! xs. P\"\n```\n``` 116 \"_Let (_binds b bs) e\" == \"_Let b (_Let bs e)\"\n```\n``` 117 \"let x = a in e\" == \"Let a (%x. e)\"\n```\n``` 118\n```\n``` 119 syntax (\"\" output)\n```\n``` 120 \"op =\" :: ['a, 'a] => bool (\"(_ =/ _)\" [51, 51] 50)\n```\n``` 121 \"op ~=\" :: ['a, 'a] => bool (\"(_ ~=/ _)\" [51, 51] 50)\n```\n``` 122\n```\n``` 123 syntax (symbols)\n```\n``` 124 Not :: bool => bool (\"\\\\<not> _\" 40)\n```\n``` 125 \"op &\" :: [bool, bool] => bool (infixr \"\\\\<and>\" 35)\n```\n``` 126 \"op |\" :: [bool, bool] => bool (infixr \"\\\\<or>\" 30)\n```\n``` 127 \"op -->\" :: [bool, bool] => bool (infixr \"\\\\<midarrow>\\\\<rightarrow>\" 25)\n```\n``` 128 \"op o\" :: ['b => 'c, 'a => 'b, 'a] => 'c (infixl \"\\\\<circ>\" 55)\n```\n``` 129 \"op ~=\" :: ['a, 'a] => bool (infixl \"\\\\<noteq>\" 50)\n```\n``` 130 \"@Eps\" :: [pttrn, bool] => 'a (\"(3\\\\<epsilon>_./ _)\" [0, 10] 10)\n```\n``` 131 \"! \" :: [idts, bool] => bool (\"(3\\\\<forall>_./ _)\" [0, 10] 10)\n```\n``` 132 \"? \" :: [idts, bool] => bool (\"(3\\\\<exists>_./ _)\" [0, 10] 10)\n```\n``` 133 \"?! \" :: [idts, bool] => bool (\"(3\\\\<exists>!_./ _)\" [0, 10] 10)\n```\n``` 134 \"@case1\" :: ['a, 'b] => case_syn (\"(2_ \\\\<Rightarrow>/ _)\" 10)\n```\n``` 135 (*\"@case2\" :: [case_syn, cases_syn] => cases_syn (\"_/ \\\\<orelse> _\")*)\n```\n``` 136\n```\n``` 137 syntax (symbols output)\n```\n``` 138 \"op ~=\" :: ['a, 'a] => bool (\"(_ \\\\<noteq>/ _)\" [51, 51] 50)\n```\n``` 139 \"*All\" :: [idts, bool] => bool (\"(3\\\\<forall>_./ _)\" [0, 10] 10)\n```\n``` 140 \"*Ex\" :: [idts, bool] => bool (\"(3\\\\<exists>_./ _)\" [0, 10] 10)\n```\n``` 141 \"*Ex1\" :: [idts, bool] => bool (\"(3\\\\<exists>!_./ _)\" [0, 10] 10)\n```\n``` 142\n```\n``` 143 syntax (xsymbols)\n```\n``` 144 \"op -->\" :: [bool, bool] => bool (infixr \"\\\\<longrightarrow>\" 25)\n```\n``` 145\n```\n``` 146 syntax (HTML output)\n```\n``` 147 Not :: bool => bool (\"\\\\<not> _\" 40)\n```\n``` 148\n```\n``` 149\n```\n``` 150 (** Rules and definitions **)\n```\n``` 151\n```\n``` 152 local\n```\n``` 153\n```\n``` 154 rules\n```\n``` 155\n```\n``` 156 eq_reflection \"(x=y) ==> (x==y)\"\n```\n``` 157\n```\n``` 158 (* Basic Rules *)\n```\n``` 159\n```\n``` 160 refl \"t = (t::'a)\"\n```\n``` 161 subst \"[| s = t; P(s) |] ==> P(t::'a)\"\n```\n``` 162\n```\n``` 163 (*Extensionality is built into the meta-logic, and this rule expresses\n```\n``` 164 a related property. It is an eta-expanded version of the traditional\n```\n``` 165 rule, and similar to the ABS rule of HOL.*)\n```\n``` 166 ext \"(!!x::'a. (f x ::'b) = g x) ==> (%x. f x) = (%x. g x)\"\n```\n``` 167\n```\n``` 168 selectI \"P (x::'a) ==> P (@x. P x)\"\n```\n``` 169\n```\n``` 170 impI \"(P ==> Q) ==> P-->Q\"\n```\n``` 171 mp \"[| P-->Q; P |] ==> Q\"\n```\n``` 172\n```\n``` 173 defs\n```\n``` 174\n```\n``` 175 True_def \"True == ((%x::bool. x) = (%x. x))\"\n```\n``` 176 All_def \"All(P) == (P = (%x. True))\"\n```\n``` 177 Ex_def \"Ex(P) == P(@x. P(x))\"\n```\n``` 178 False_def \"False == (!P. P)\"\n```\n``` 179 not_def \"~ P == P-->False\"\n```\n``` 180 and_def \"P & Q == !R. (P-->Q-->R) --> R\"\n```\n``` 181 or_def \"P | Q == !R. (P-->R) --> (Q-->R) --> R\"\n```\n``` 182 Ex1_def \"Ex1(P) == ? x. P(x) & (! y. P(y) --> y=x)\"\n```\n``` 183\n```\n``` 184 rules\n```\n``` 185 (* Axioms *)\n```\n``` 186\n```\n``` 187 iff \"(P-->Q) --> (Q-->P) --> (P=Q)\"\n```\n``` 188 True_or_False \"(P=True) | (P=False)\"\n```\n``` 189\n```\n``` 190 defs\n```\n``` 191 (*misc definitions*)\n```\n``` 192 Let_def \"Let s f == f(s)\"\n```\n``` 193 if_def \"If P x y == @z::'a. (P=True --> z=x) & (P=False --> z=y)\"\n```\n``` 194\n```\n``` 195 (*arbitrary is completely unspecified, but is made to appear as a\n```\n``` 196 definition syntactically*)\n```\n``` 197 arbitrary_def \"False ==> arbitrary == (@x. False)\"\n```\n``` 198\n```\n``` 199\n```\n``` 200\n```\n``` 201 (** initial HOL theory setup **)\n```\n``` 202\n```\n``` 203 setup Simplifier.setup\n```\n``` 204 setup ClasetThyData.setup\n```\n``` 205\n```\n``` 206\n```\n``` 207 end\n```\n``` 208\n```\n``` 209\n```\n``` 210 ML\n```\n``` 211\n```\n``` 212\n```\n``` 213 (** Choice between the HOL and Isabelle style of quantifiers **)\n```\n``` 214\n```\n``` 215 val HOL_quantifiers = ref true;\n```\n``` 216\n```\n``` 217 fun alt_ast_tr' (name, alt_name) =\n```\n``` 218 let\n```\n``` 219 fun ast_tr' (*name*) args =\n```\n``` 220 if ! HOL_quantifiers then raise Match\n```\n``` 221 else Syntax.mk_appl (Syntax.Constant alt_name) args;\n```\n``` 222 in\n```\n``` 223 (name, ast_tr')\n```\n``` 224 end;\n```\n``` 225\n```\n``` 226\n```\n``` 227 val print_ast_translation =\n```\n``` 228 map alt_ast_tr' [(\"! \", \"*All\"), (\"? \", \"*Ex\"), (\"?! \", \"*Ex1\")];\n```" ]
[ null ]
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https://ebrary.net/113969/communication/three_beam_interference
[ "", null, "# Three-Beam Interference\n\nGenerally speaking, the principle of three-beam interference is based on the interference of three reflective surfaces. Figure 1.5 shows one of the typical structures of the three-beam interference based on graded- index multimode fiber (GI-MMF) . It is fabricated by cascading an air gap and a short section of GI-MMF to a singlemode fiber (SMF). The air gap can be formed by fusion splicing the SMF with a chemically etched micronotch on the GI-MMF end. The light propagation follows the sinusoidal path as the index profile is parabolic.\n\nThe ray-transfer-matrix (RTM) theory is used to describe the principle. A Gaussian beam can be expressed as", null, "where A0 remains constant for energy conservation. k = 2n/X is the wave number X with the free-space wavelength. Ф = nkz denotes the phase shift of the light beam as it propagates and n is the refractive", null, "index of the medium light propagating. The complex beam parameter, q, is given by", null, "where p and Ю are the radius of the curvature and the beam radius of the Gaussian beam, respectively. Suppose M is the transfer matrix from input plane to output plane, the transformation of the complex beam parameter is given by", null, "where A = M(1,1), B = M(1,2), C = M(2,1), and D = M(2,2), and q and q are the complex beam parameters at the input and output planes, respectively.\n\nThe electrical amplitude of the incident beam at location 0 is E0(r) and corresponding complex parameter is q0 = insp&2S /l. ns is the reflective index of the SMF core and is the beam radius of the SMF. The reflectance of surface Rj is determined by the Fresnel equation, Rj = (ns n0)2/(ns + n0)2, and n0 is the reflective index of the ambient medium, which is approximately 1 in the air. The electrical amplitude of the reflective beam at surface I is E1( r) = -y/R\"E0 (r). The reflectance of the etched micronotch (surface II) is represented by", null, "where and ЮПг are the beam radii of the incident light and the reflected light by surface II, respectively. R(r) is the reflectance given by the radial distribution, R(r) = [n(r) - n0]2/[n(r) + n0]2, where n(r) is the refractive index profile of the GI-MMF core, which can be given by", null, "where n1 is the maximum index at r = 0, a is the radius of the GI-MMF core, and g is a factor that determines the index profile of the core. The ABCD matrices corresponding to EIIr and Em are given by Mm = M12M01 and MIIr = M2M12M01, where Mj is the matrix describing the transformation of the complex beam parameters between locations i and j. The elementary matrices are given by", null, "where p1 is the radius of curvature of the etched micronotch on the GI-MMF end and L0 is the effective cavity length of the air gap, which is smaller than the distance between the SMF end and the bottom of the etched micronotch. The transformation of the complex beam parameters is determined by Equation 1.13. The electrical amplitude of the light beam reflected by surface II is EII(r) = R{IEI,I(r), where RII = Т— (1 — Aj) Rn. TI is the transmittance of surface I and AI is the propagation losses in the air gap. Фп, the additional phase in En, is given by Фп = 2n0kL0. The ABCD matrix corresponding\n\nto En is presented as Mn = M10M2rM2M12M01 with M2V = M12 and\n\nГ1 о ]\n\nMin = ~ / . The effective reflectance of surface II can be\n\n0 no /ns\n\nexpressed as", null, "where 2as is the mode field diameter of the SMF.\n\nAnalogizing to the reflection Rn, the reflectance of surface III is given by", null, "The electrical amplitude of the light beam reflected by surface 111 is Em{r) = ylR{nЕ{ц(г) with Ащ = 7]2Tn(1 - An)Am. 7} is the transmittance of surface I and AI is the propagation losses in the air gap. Фш, the additional phase in Em, is given by Фш = 2(n0L0 + n1L). The ABCD matrices corresponding to EIffi, EIIIr, and EIII are presented as Mmi = M34M23M12M01> MIIIr = M44,M34M23M12M01> and\n\nMIII = My0’M2YM32 M33/M23M12M01, respectively, with", null, "Considering the coupling losses of the light beam reflected from III into the SMF, the effective reflectance of surface III can be expressed as", null, "R тф is mainly dependent on the coupling coefficient of the light beam into the SMF as the GI-MMF length changes. In general, RnIff is smaller than RIII, because of the coupling losses and propagation losses. The reflective signal of the three-beam interference is given by", null, "By using the effective reflectance of the three-beam interferometer, the above RTM theory can be simplified. For the case of n0 < n1, the normalized intensity of the three-beam interference can be expressed as", null, "It is well known that the fringe contrast of the two-beam interference becomes maximum when the reflectance of the surfaces is equal, which is a strict constraint. The corresponding constraint condition on the effective reflectance of the three surfaces for the three-beam interference can be deduced by Equation 1.10 and is given by", null, "Unlike the two-beam interference, the constraint condition for the three-beam interference to obtain the optimal fringe contrast is an inequality; that is, the requirement on the reflectance is a relatively wide range rather than a decided value. This makes it easier for sensing based on three-beam interference to obtain high performance than the conventional two-beam interference. Figure 1.6 shows the reflective spectra of the SMF end (black), the air gap (light gray), and the three-beam interferometer introduced above with GI-MMF length of 515 pm (gray) .", null, "Figure 1.6 Reflective spectra of two- and three-beam interferences.", null, "" ]
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https://chemistry.stackexchange.com/questions/28646/how-to-write-equations-for-dissociation-of-ionic-compounds-in-water
[ "# How to write equations for dissociation of ionic compounds in water?\n\nPlease write a dissociation-in-water equation for the compounds $\\ce{BeI2}$ and $\\ce{LiI}$. Make sure to add the states of matter after each compound.\n\nCurrently, for $\\ce{BeI2}$ I have the equation $$\\ce{BeI2 (s) -> Be^2+ (s) + I2^2- (g)}.$$ I have yet to attempt the second one.\n\n## 2 Answers\n\nYour ionic charges are not correct for iodine. Looking at your attempt.\n\nAs $\\ce{Be}$ is in group 2, the ionic charge for beryllium ion is fine, but iodine is in group 17, so its ion is $\\ce{I-}$. When the ions dissociate, they become aqueous or (aq) as the state of matter.\n\nThen the ionic charges need to balance, thus:\n\n$$\\ce{BeI2 (s) -> Be^2+ (aq) + 2I- (aq)}$$\n\nTo balance the ionic charges in this example, you need 2 $\\ce{I-}$ (iodine ions) to balance the $\\ce{Be^{2+}}$ (beryllium ion).\n\nA similar example (and further explanations) are provided on the UC Davis ChemWiki page Unique Features of Aqueous Solutions (including an example of the dissolution of $\\ce{MgCl2}$ - another compound with group 2 and 17 elements).\n\nSo,\n\n• determine the group, hence ionic charge of each dissociated ion\n• balance these charges\n• state that the dissociated ions are aqueous\n\nNow, use the process to determine the dissociation of $\\ce{LiI}$\n\nBy definition, a dissociation-in-water reaction results in aqueous, independent ions, not negatively charged diatomic molecules like the $\\ce{I_2^2-}$ in the original question (not to mention that diatomic molecules are generally not negatively charged).\n\nThe answer should result in solely single-atom ions, as so:\n\n$$\\ce{BeI2 (s) -> Be^2+ (aq) + 2I- (aq)}$$\n\nAnd using this same logic for $\\ce{LiI}$:\n\n$$\\ce{LiI (s) -> Li+ (aq) + I- (aq)}$$" ]
[ null ]
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https://link.springer.com/article/10.3103%2FS0146411619060038
[ "# Neural Networks Based Prediction Model for Vessel Track Control\n\n• 18 Accesses\n\n### Abstract\n\nThe problem of neural networks implementation for the construction of a predictive model for vessel track control was studied. It has been shown that the vessel track control problem may be considered as an approximation task, and neural networks may be implemented as universal approximating tools. The general structure of the prediction model, based on neural networks, has been developed. The model consists of several two-layered feedforward neural networks, which architectures satisfy the conditions of universal approximation properties. The analysis of the functions of the different neural networks in the prediction model has been performed. The network predicting WGS-84 geodetic latitude as a part of the predictive model has been constructed, trained and validated by using MATLAB software. The validation results show the good prediction precision of the net.\n\nThis is a preview of subscription content, log in to check access.\n\n## Access options\n\nUS\\$ 39.95\n\nPrice includes VAT for USA\n\n1. 1.\n\nFor clarity, let suppose that the route line is a number of waypoints on the Earth ellipsoid, connected with the segments of geodetic lines.\n\n## REFERENCES\n\n1. 1\n\nChen, S., et al., Approximating explicit model predictive control using constrained neural networks, Proc. of the American Control Conf., Milwauke, 2018, pp. 1520–1527.\n\n2. 2\n\nCybenko, G., Approximation by superpositions of a sigmoidal function, Math. Control Signals Syst., 1989, vol. 2, pp. 303–314.\n\n3. 3\n\nHaykin, S., Neural Networks and Learning Machines, New York: Prentice Hall, 2009.\n\n4. 4\n\nKainen, P.C., Kurkova, V., and Sanguineti, M., Approximating multivariable functions by feedforward neural nets, in Handbook on Neural Information Processing, New York: Springer, 2013, ch. 5.\n\n5. 5\n\nLeshno, M., et al., Multilayer feedforward networks with a nonpolynomial activation function can approximate any function, Neural Networks, 1993, vol. 6, pp. 861–867.\n\n6. 6\n\nLawrynczuk, M., Training of neural models for predictive control, Neurocomputing, 2010, vol. 73, pp. 1332–1343.\n\n7. 7\n\nPatan, K., Two stage neural network modeling for robust model predictive control, ISA Trans., 2018, pp. 56–65.\n\n8. 8\n\nRanković, V., et al., Neural network model predictive control of nonlinear systems using genetic algorithms, Int. J. Comput. Commun., 2012, vol. 7, no. 3, pp. 540–549.\n\n9. 9\n\nVincent, A.A. and Hassapis, G., Adaptive predictive control using recurrent neural network identification, Proc. of 17th Mediterranean Conf. on Control & Automotion, Thessaloniki, 2009, pp. 61–66.\n\n10. 10\n\nSoloway, D. and Pamela, J.H., Neural generalized predictive control. A Newton-Raphson implementation, Proceedings of the 1996 IEEE International Symp. on Intelligent Control, Dearborn, 1996, pp. 277–282.\n\n11. 11\n\nReese, B. and Collins, E., A graph search and neural network approach to adaptive nonlinear model predictive control, Eng. Appl. Artif. Intell., 2016, vol. 55, pp. 250–268.\n\n12. 12\n\nYan, Z. and Wang, J., Model predictive control for tracking of underactuated vessels based on recurrent neural networks, IEEE J. Oceanic Eng., 2012, vol. 37, no. 4, pp. 717–726.\n\n13. 13\n\nFossen, T.I., Handbook of Marine Craft Hydrodynamics and Motion Control, Chichester: John Wiley & Sons, 2011.\n\n14. 14\n\nFaltinsen, O.M., Sea Loads on Ships and Offshore Structures, Cambridge: Cambridge University Press, 1990.\n\n15. 15\n\nNguyen, D. and Widrow, B., Improving the learning speed of 2-layer neural networks by choosing initial values of the adaptive weights, Proc. of 1990 IJCNN International Joint Conference on Neural Networks, San Diego, 1990, vol. 3, pp. 21–26.\n\n16. 16\n\nForesee, F.D. and Hagan, M.T., Gauss-Newton approximation to Bayesian learning, Proc. of International Conference on Neural Networks (ICNN’97), Houston, 1997, vol. 3, pp. 1930–1935." ]
[ null ]
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https://au.mathworks.com/matlabcentral/answers/471730-change-the-iteration-in-for-loop?s_tid=prof_contriblnk
[ "# change the iteration in for loop\n\n41 views (last 30 days)\nNA on 15 Jul 2019\nCommented: Andy on 16 Jul 2019\nI have For Loop that calculate c. Sometimes c become nan or big number.\nWhen I get big amount or Nan for c, I want to repeat that iteration again.\nI used this code but does not repeat that iteration.\nfor i=1:10\n% calculate c\nif isnan(c(i))==1 | c(i)>0.009\ni=i-1 % again repeat the iteration\nend\nend\n\nKALYAN ACHARJYA on 15 Jul 2019\nI think '|' and '||' does not have a difference.\nwhen it goes to if part, c become zero on that iteration.\nit means that does not repeat that iteration.\nfor example in i=2, c(2)=nan. when I use this code i=1, but in for loop it starts from 3.\nBig number in my case more than 0.009\nKALYAN ACHARJYA on 15 Jul 2019\nm=m-1; % matlab gives me this detail\nit appears that the index value of the indicated for loop changes inside the loop body. Often, this happens when loops nest and an inner loop reuses the name of an outer loop index. Because MATLAB resets the loop index to the next value when it returns to the top of the outer loop, it ignores any changes that took place within a nested loop. If the code does not to refer to the outer loop index value after the inner loop changes it, no problem results. However if you attempt to use the outer loop index value after the inner loop completes, it is likely to result in a bug.\n\nAndy on 15 Jul 2019\nI don't think the For loop is what you need.\ni =1;\nwhile i<11\n% calculate c\nif isnan(c(i))==0 & c(i)<=0.009\ni=i+1;\nend\nend\n\nNA on 15 Jul 2019\nIf i=5; nan happens,\ni changes to 6, but c(5) is still nan.\nI think i=i+1; does not work\nAndy on 16 Jul 2019\nI made up this code, adding the else just to change the value and it works fine.\ni=1;\nc=[ .001 12 .003 .004 nan .006 .007 .008 nan .005];\nwhile i<11\n%calculate c\nif isnan(c(i))==0 & c(i)<=.009\ni=i+1;\nelse\nc(i)\nc(i)=0;\nend\nend\nc\n\nKALYAN ACHARJYA on 15 Jul 2019\nEdited: KALYAN ACHARJYA on 15 Jul 2019\n# Experts need your suggestions here\nfor i=1:10\nc(i)=...% do\nwhile isnan(c(i)) || c(i)>0.009\nc(i)=..\nend\nend\nI know multiple loop is messed here, just try to get way out. Hope I undestand the question.\nOne Note: Without changing i, is there any possibilty to change the C(i) in next or next iterartions within while loop, so that once it fail, it exit from while loop?\nThe code runs within the while loop, without changing i ultill C(i) satisfy any one conditions-\n1. C(i) is NaN\n2. C(i)>0.009\n\nNA on 15 Jul 2019\nfor i=1:10\nc(i)=...% do\nwhile isnan(c(i)) || c(i)>0.009\nc(i)=[]; % ignore calculated c(i)\nend\nend\nresult of c(i)\n0.0038 0.0036 0.0029 0 0.0019 0.0011 0.0018\n0.0011 0.0305\nwhen i=4, nan happens so I want to repeat i=4 to get some result.\nKALYAN ACHARJYA on 15 Jul 2019\nwhen i=4, nan happens so I want to repeat i=4 to get some result\nWhen NaN appears, it enter to within while loop, while loop doesnot change the i value. When c(i)=Nan at i=4, it keep running withing while loop, that why I asking, is there any possibility to change C(i) result without changing i, so that it exit from while loop?" ]
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https://safecurves.cr.yp.to/proof/110244005101.html
[ "Primality proof for n = 110244005101:\n\nTake b = 2.\n\nb^(n-1) mod n = 1.\n\n13610371 is prime.\nb^((n-1)/13610371)-1 mod n = 90519940136, which is a unit, inverse 61327983725.\n\n(13610371) divides n-1.\n\n(13610371)^2 > n.\n\nn is prime by Pocklington's theorem." ]
[ null ]
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https://www.devzoneoriginal.com/2020/08/convert-h5-keras-model-into-tflite-model.html
[ "### [SOLVED] Convert .h5 Keras model into the .tflite model\n\nDue to TensorFlow versions compatibility issues, you might face issues in converting .h5 Keras model into the .tflite model. I have found a workaround for this and going to share it with you.\n\nStep 1: Uninstall the current version of TensorFlow, by typing below command:\n\npip uninstall tensorflow\n\nIf you are using Google Colab, use\n\n!pip uninstall tensorflow\n\nStep 2: Install TensorFlow version 1.12, by typing below command:\n\npip install tensorflow==1.12\n\nIf you are using Google Colab, use\n\n!pip install tensorflow==1.12\n\nStep 3: Add the import statement\n\nfrom tensorflow.contrib import lite\n\nStep 4: Code snippet for model conversion\n\nconverter = lite.TFLiteConverter.from_keras_model_file('/content/my_model.h5')\n\ntfmodel = converter.convert()\n\nopen(\"model.tflite\", \"wb\").write(tfmodel)\n\nAnd that's it. You can now use the generated model.tflite file to perform the inferences." ]
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https://mathhelpboards.com/threads/problem-of-the-week-57-july-1st-2013.5407/
[ "# Problem of the Week #57 - July 1st, 2013\n\nStatus\nNot open for further replies.\n\n#### Chris L T521\n\n##### Well-known member\nStaff member\nHere's this week's problem.\n\n-----\n\nProblem: Let $f$ be continuous on a domain $\\Omega\\subseteq\\mathbb{C}$ and assume that $\\displaystyle\\int_Tf(z)\\,dz=0$ for each triangle $T\\subset\\Omega$. Show that $f$ is holomorphic.\n\n-----\n\n#### Chris L T521\n\n##### Well-known member\nStaff member\nNo one answered this week's question. You can find my solution below.\n\nProof: Let $\\mathbb{D}(p,r)=D\\subset\\Omega$ be a disc. Define a function $F: D\\rightarrow\\mathbb{C}$ by $F(z)=\\int_{[p,z]} f(\\zeta)\\,d\\zeta\\qquad z\\in D.$\nFix some $z_0\\in\\Omega$. Then if $T=[p,z]\\cup [z,z_0]\\cup [z_0,p]$ (i.e. the triangle $T$ is made of those three line segments in $D$), we see that\n\\begin{aligned} \\int_T f(\\zeta)\\,d\\zeta=0 &\\implies \\int_{[p,z]} f(\\zeta)\\,d\\zeta + \\int_{[z,z_0]} f(\\zeta)\\,d\\zeta + \\int_{[z_0,p]} f(\\zeta)\\,d\\zeta = 0\\\\ &\\implies \\int_{[p,z]}f(\\zeta)\\,d\\zeta - \\int_{[p,z_0]} f(\\zeta)\\,d\\zeta) = \\int_{[z_0,z]} f(\\zeta)\\,d\\zeta.\\end{aligned}\nHence,\n$\\frac{F(z)-F(z_0)}{z-z_0}-f(z_0) = \\frac{1}{z-z_0} \\int_{[z_0,z]}\\left( f(\\zeta)-f(z_0)\\right)\\,d\\zeta.$\nSince $f$ is continuous at $z_0$, then for each $\\varepsilon>0$ there is a $\\delta>0$ such that $|f(\\zeta)-f(z_0)|<\\varepsilon$ whenever $|\\zeta-z_0|<\\delta$. This now implies that\n\\begin{aligned} \\left|\\frac{F(z)-F(z_0)}{z-z_0} - f(z_0)\\right| &= \\left|\\frac{1}{z-z_0}\\int_{[z_0,z]} \\left( f(\\zeta)-f(z_0)\\right)\\,d\\zeta\\right| \\\\ &\\leq \\left|\\frac{1}{z-z_0}\\right| \\int_{[z_0,z]} |f(\\zeta)-f(z_0)|\\,d\\zeta\\\\ &\\leq \\varepsilon \\left|\\frac{1}{z-z_0}\\right| \\int_{[z_0,z]}\\,d\\zeta\\\\ &= \\varepsilon\\end{aligned}\nHence,\n$\\left|\\frac{F(z)-F(z_0)}{z-z_0} - f(z_0)\\right| \\leq \\varepsilon\\qquad \\text{if |z-z_0|<\\delta}$\nIt now follows that $F^{\\prime}(z_0)$ exists and equals $f(z_0)$. Since $z_0\\in D$ was chosen arbitrarily, it follows that $F^{\\prime}=f$ (i.e. $f$ has a primitive); furthermore, we have that $F$ is holomorphic in $D$. Since the derivative of a holomorphic function is holomorphic, it now follows that $f$ is holomorphic in $D$. Furthemore, since this is true for every disk $D$ contained in $\\Omega$, it must follow that $f$ is holomorphic in $\\Omega$.$\\hspace{.25in}\\blacksquare$\n\nStatus\nNot open for further replies." ]
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https://www.assignmentexpert.com/homework-answers/physics/mechanics-relativity/question-109472
[ "# Answer to Question #109472 in Mechanics | Relativity for let me know\n\nQuestion #109472\n1\n2020-04-15T10:43:28-0400\n\nAs per the given question,\n\nMass of the block A(m)=4kg\n\nMass of the cylinder (M)=6kg\n\nLet the tension in the string is T and the acceleration of the block is a, cylinder is rolling without slipping.", null, "\"mg-T=ma\"\n\n\"4g-T=4a-----(i)\"\n\nNow, taking torque about the point of contact,\n\n\"I\\\\alpha = TR\"\n\n\"T=\\\\dfrac{I\\\\alpha}{R}=(MR^2+\\\\dfrac{MR^2}{2})\\\\dfrac{\\\\alpha}{R}=\\\\dfrac{3MR^2\\\\times \\\\alpha}{2R}=\\\\dfrac{3Ma}{2}----(ii)\"\n\nFrom equation (i) and (ii)\n\n\"\\\\Rightarrow 4g+\\\\dfrac{3Ma}{2}=4a\"\n\n\"\\\\Rightarrow 4a+\\\\dfrac{3\\\\times 6a}{2}=4g\"\n\n\"\\\\Rightarrow 4a+9a=4g\"\n\n\"\\\\Rightarrow 13a=4g\"\n\n\"\\\\Rightarrow a=\\\\dfrac{4g}{13}=3.01 m\\/sec^2\"\n\nTension in the string, \"T=\\\\dfrac{3Ma}{2}=\\\\dfrac{3\\\\times 6\\\\times 3.01}{2}27.09N\"\n\nii) Speed of the cylinder after t=5 sec\n\n\"v=u+at\"\n\n\"v=0+3.01\\\\times 5= 15.05 m\\/sec\"\n\nangular velocity at t=5,\n\n\"\\\\omega =\\\\dfrac{v}{R}=\\\\dfrac{15.05}{0.2}=75.25 rev\\/sec\"\n\nHence kinetic energy of the cylinder =\"\\\\dfrac{MR^2}{2}+\\\\dfrac{I\\\\omega^2}{2}=\\\\dfrac{6\\\\times 15.05^2}{2}+\\\\dfrac{3\\\\times6\\\\times 0.2^2\\\\times 75.25^2}{2}\"\n\n\"=679.50 +2038.52=2718.02J\"\n\nNeed a fast expert's response?\n\nSubmit order\n\nand get a quick answer at the best price\n\nfor any assignment or question with DETAILED EXPLANATIONS!" ]
[ null, "https://www.assignmentexpert.com/image", null ]
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https://www.numbers.education/10070.html
[ "Is 10070 a prime number? What are the divisors of 10070?\n\n## Parity of 10 070\n\n10 070 is an even number, because it is evenly divisible by 2: 10 070 / 2 = 5 035.\n\nFind out more:\n\n## Is 10 070 a perfect square number?\n\nA number is a perfect square (or a square number) if its square root is an integer; that is to say, it is the product of an integer with itself. Here, the square root of 10 070 is about 100.349.\n\nThus, the square root of 10 070 is not an integer, and therefore 10 070 is not a square number.\n\n## What is the square number of 10 070?\n\nThe square of a number (here 10 070) is the result of the product of this number (10 070) by itself (i.e., 10 070 × 10 070); the square of 10 070 is sometimes called \"raising 10 070 to the power 2\", or \"10 070 squared\".\n\nThe square of 10 070 is 101 404 900 because 10 070 × 10 070 = 10 0702 = 101 404 900.\n\nAs a consequence, 10 070 is the square root of 101 404 900.\n\n## Number of digits of 10 070\n\n10 070 is a number with 5 digits.\n\n## What are the multiples of 10 070?\n\nThe multiples of 10 070 are all integers evenly divisible by 10 070, that is all numbers such that the remainder of the division by 10 070 is zero. There are infinitely many multiples of 10 070. The smallest multiples of 10 070 are:\n\n## Numbers near 10 070\n\n### Nearest numbers from 10 070\n\nFind out whether some integer is a prime number" ]
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https://www.physicsforums.com/threads/kinematics-in-2d-captain-of-a-plane-question.341827/
[ "# Kinematics in 2D - Captain of a Plane Question\n\ncrono_\nThis is an assignment question, so I'm not looking for the answer but perhaps just a nudge in the right direction.\n\n## Homework Statement\n\nThe captain of a plane wishes to proceed due west. The cruising speed of the plane is 260 m/s relative to the air. A weather report indicates that a 45.0-m/s wind is blowing from the south to the north. In what direction, measured with respect to due west, should the pilot head the plane?\n\n## Homework Equations\n\nThis sounds like it's a right triangle, so I think Pythagorean theorem would be used.\n\nc2 = a2 + b2\n\n## The Attempt at a Solution\n\nI was attempting to draw the problem on a standard x, y coordinate system but wasn't sure how to lay it out.\n\nThe wind seems pretty obvious as it's blowing from the south to north, so I was thinking about using a vector arrow pointing in the + y direction. The cruising speed of the plane however seems a little trickier. Would I consider this to be the hypotenuse? Or just one of the other two sides of the right triangle?\n\nOnce the two values have been drawn I imagine you would use Pythagorean to find the other side and then calculate the angle." ]
[ null ]
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http://degreeoffer.info/true-sine-inverter/true-sine-inverter-mkp1000-241r-hot-sale-off-grid-pure-sine-wave-power-inverter-1000-watt-true-free-sine-wave-inverter-schematic/
[ "True Sine Inverter Mkp1000 241r Hot Sale Off Grid Pure Sine Wave Power Inverter 1000 Watt True Free Sine Wave Inverter Schematic", null, "true sine inverter mkp1000 241r hot sale off grid pure sine wave power inverter 1000 watt true free sine wave inverter schematic.\n\ntrue sine wave inverter reviews,low harmonic distortion true sine inverter,what is a true sine inverter,true sine inverter reviews,no distrotion true sine inverter,free sine wave inverter schematic,true sine wave inverter charger,best sine wave inverter generators,true sine wave inverter schematic,true sine inverter,true sine inverter charger." ]
[ null, "http://degreeoffer.info/data/true-sine-inverter/images/true-sine-inverter-mkp1000-241r-hot-sale-off-grid-pure-sine-wave-power-inverter-1000-watt-true-free-sine-wave-inverter-schematic.jpg", null ]
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https://wordpress.org/support/topic/does-not-work-function/
[ "Support » Plugin: Calculated Fields Form » does not work “function”\n\n• Hi,\nEnter a formula into a calculated field. But only the first condition works. For the rest, 0 is displayed. Where is the mistake?\n\n``````(function(){\nvar perem;\nif(1 <= fieldname9, fieldname9 <= 20) {\nperem = 1;\n}\nelse if (21<= fieldname9, fieldname9 <= 40) {\nperem = 0,7;\n}\nelse if (41<= fieldname9, fieldname9 <= 50) {\nperem = 0,5;\n}\nelse if (51<= fieldname9, fieldname9 <= 10000) {\nperem = 0,3;\n}\n\nreturn Math.round(perem*fieldname9*fieldname11);\n})()``````\nViewing 2 replies - 1 through 2 (of 2 total)\n• Hello @anka3000\n\nThank you very much for using our plugin. Please, include the “and” (&&) operator in the conditional statements. Furthermore, in javascript, the decimal symbol in decimal numbers is the point:\n\n``````(function(){\nvar perem;\nif(1 <= fieldname9 && fieldname9 <= 20) {\nperem = 1;\n}\nelse if (21<= fieldname9 && fieldname9 <= 40) {\nperem = 0.7;\n}\nelse if (41<= fieldname9 && fieldname9 <= 50) {\nperem = 0.5;\n}\nelse if (51<= fieldname9 && fieldname9 <= 10000) {\nperem = 0.3;\n}\n\nreturn ROUND(perem*fieldname9*fieldname11);\n})()``````\n\nBest regards.\n\nThank you! You are the best!!!\n\nViewing 2 replies - 1 through 2 (of 2 total)\n• You must be logged in to reply to this topic." ]
[ null ]
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https://mathfilefoldergames.com/2014/11/03/lets-find-the-10s-free/
[ "# Let’s Find the 10s {FREE}", null, "You can do this activity with your students at the board or on worksheets.\n\nProvide them with a really long horizontal sum. Something like this:\n\n3 + 11 + 7 + 5 + 5 + 4 + 9 + 1 + 2 + 13 + 8 =\n\nThe idea here is to group pairs of numbers into tens or multiples of ten so that they can quickly and easily find the sum using mental math.\n\nHere there are four pairs that yield 10: (3 + 7); (5 + 5); (9 +1); (2 + 8) so that total is 40 plus 11 + 4 + 13 = 28\n\nSo the total sum is 40 + 28 = 68\n\nStudents who are agile with this will come up with other ways to get the same answer.\n\nFor example: (7 + 13); (11 + 9) = 40\nThen, (5 + 5); (3 + 4 + 1 + 2) = 20", null, "" ]
[ null, "https://i2.wp.com/www.mathfilefoldergames.com/wp-content/uploads/2014/11/Text1-1024x815.png", null, "https://i2.wp.com/www.mathfilefoldergames.com/wp-content/uploads/2014/11/MentalMathBundle-ClickHere-1024x768.png", null ]
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https://shiji.info/tag/easy/
[ "# 1089. Duplicate Zeros\n\nGiven a fixed length array arr of integers, duplicate each occurrence of zero, shifting the remaining elements to the right.\n\nNote that elements beyond the length of the original array are not written.\n\nDo the above modifications to the input array in place, do not return anything from your function.\n\nExample 1:\n\nInput: [1,0,2,3,0,4,5,0]\nOutput: null\nExplanation: After calling your function, the input array is modified to: [1,0,0,2,3,0,0,4]\n\nExample 2:\n\nInput: [1,2,3]\nOutput: null\nExplanation: After calling your function, the input array is modified to: [1,2,3]\n\nNote:\n\n``````1 <= arr.length <= 10000\n0 <= arr[i] <= 9\n``````\n\nSolution:\n\n```class Solution {\npublic void duplicateZeros(int[] arr) {\nint count = 0;\nfor (int i = 0; i < arr.length; i ++){\nif (arr[i] == 0){\ncount ++;\n}\n}\n\nint offset = count;\nfor (int i = arr.length - 1; i >= 0; i --){\nif (i + offset < arr.length){\narr[i+offset] = arr[i];\n}\nif (arr[i] == 0){\noffset --;\nif (i + offset < arr.length){\narr[i+offset] = arr[i];\n}\n}\n}\n}\n}```\n\n# 88. Merge Sorted Array\n\nGiven two sorted integer arrays nums1 and nums2, merge nums2 into nums1 as one sorted array.\n\nNote:\n\nThe number of elements initialized in nums1 and nums2 are m and n respectively.\nYou may assume that nums1 has enough space (size that is equal to m + n) to hold additional elements from nums2.\n\nExample:\n\nInput:\nnums1 = [1,2,3,0,0,0], m = 3\nnums2 = [2,5,6], n = 3\n\nOutput: [1,2,2,3,5,6]\n\nConstraints:\n\n-10^9 <= nums1[i], nums2[i] <= 10^9 nums1.length == m + n nums2.length == n Solution.\n\n```class Solution {\npublic void merge(int[] nums1, int m, int[] nums2, int n) {\n// index for m\nint im = m -1;\n// index for n\nint in = n -1;\n// index for result\nint ir = m + n -1;\n\nfor (;ir >= 0; ir --){\n\nif (im < 0){\n// If one of the array is all cleared, use the other\nnums1[ir] = nums2[in--];\n} else if (in < 0){\nnums1[ir] = nums1[im--];\n} else if (nums1[im] > nums2[in]){\n// Otherwise, pick the greater one.\nnums1[ir] = nums1[im--];\n} else {\nnums1[ir] = nums2[in--];\n}\n}\n\n}\n}\n```\n\n# 1170. Compare Strings by Frequency of the Smallest Character\n\nLet’s define a function `f(s)` over a non-empty string `s`, which calculates the frequency of the smallest character in `s`. For example, if `s = \"dcce\"` then `f(s) = 2` because the smallest character is `\"c\"` and its frequency is 2.\n\nNow, given string arrays `queries` and `words`, return an integer array `answer`, where each `answer[i]` is the number of words such that `f(queries[i])` < `f(W)`, where `W` is a word in `words`.\n\nExample 1:\n\nInput:\n\n``` queries = [\"cbd\"], words = [\"zaaaz\"]\n```\n\nOutput:\n\n``` \n```\n\nExplanation:\n\n``` On the first query we have f(\"cbd\") = 1, f(\"zaaaz\") = 3 so f(\"cbd\") < f(\"zaaaz\").\n```\n\nExample 2:\n\nInput:\n\n``` queries = [\"bbb\",\"cc\"], words = [\"a\",\"aa\",\"aaa\",\"aaaa\"]\n```\n\nOutput:\n\n``` [1,2]\n```\n\nExplanation:\n\n``` On the first query only f(\"bbb\") < f(\"aaaa\"). On the second query both f(\"aaa\") and f(\"aaaa\") are both > f(\"cc\").\n```\n\nConstraints:\n\n• `1 <= queries.length <= 2000`\n• `1 <= words.length <= 2000`\n• `1 <= queries[i].length, words[i].length <= 10`\n• `queries[i][j]``words[i][j]` are English lowercase letters.\n\n#### Straight forward solution O(q * w)\n\n```class Solution {\npublic int[] numSmallerByFrequency(String[] queries, String[] words) {\nint[] result = new int[queries.length];\nfor (int i = 0 ; i < queries.length; i ++){\nint score = s(queries[i]);\n\nfor (String w : words){\nif (s(w) > score){\nresult[i] ++;\n}\n}\n}\nreturn result;\n}\n\npublic static int s(String s) {\nchar min = 'z';\nchar[] ca = s.toCharArray();\nfor(char c : ca){\nif (c < min){\nmin = c;\n}\nif (c == 'a'){\nbreak;\n}\n}\nint r = 0;\nfor(char c : ca){\nif (c == min)\nr ++;\n}\nreturn r;\n}\n}```\n\n#### Optimized solution O (q + w)\n\nsince we have `1 <= queries[i].length, words[i].length <= 10`\n\nwhich means query will have only 10 possible values (limited), so for the result.\n\nso we can use counting.\n\n```class Solution {\npublic int[] numSmallerByFrequency(String[] queries, String[] words) {\n\nint[] freqCounter = new int;\nint[] result = new int[queries.length];\nfor (int i = 0; i < words.length; i++) {\nint freq = s(words[i]);\nfreqCounter[freq]++;\n}\n\nfor (int i = 1; i < freqCounter.length; i++) {\nfreqCounter[i] += freqCounter[i - 1];\n}\nfor (int i = 0; i < queries.length; i++) {\nint freq = s(queries[i]);\nresult[i] = words.length - freqCounter[freq];\n}\n\nreturn result;\n}\n\npublic static int s(String s) {\nchar min = 'z';\nint r = 0;\nchar[] ca = s.toCharArray();\nfor(char c : ca){\nif (c < min){\nr = 1;\nmin = c;\n} else if(c == min) {\nr ++;\n}\n\n}\nreturn r;\n}\n\n}```\n\n# 1275. Find Winner on a Tic Tac Toe Game\n\nTic-tac-toe is played by two players A and B on a 3 x 3 grid.\n\nHere are the rules of Tic-Tac-Toe:\n\n• Players take turns placing characters into empty squares (” “).\n• The first player A always places “X” characters, while the second player B always places “O” characters.\n• “X” and “O” characters are always placed into empty squares, never on filled ones.\n• The game ends when there are 3 of the same (non-empty) character filling any row, column, or diagonal.\n• The game also ends if all squares are non-empty.\n• No more moves can be played if the game is over.\n\nGiven an array `moves` where each element is another array of size 2 corresponding to the row and column of the grid where they mark their respective character in the order in which A and B play.\n\nReturn the winner of the game if it exists (A or B), in case the game ends in a draw return “Draw”, if there are still movements to play return “Pending”.\n\nYou can assume that `moves` is valid (It follows the rules of Tic-Tac-Toe), the grid is initially empty and A will play first.\n\nExample 1:\n\nInput:\n\n``` moves = [[0,0],[2,0],[1,1],[2,1],[2,2]]\n```\n\nOutput:\n\n``` \"A\"\n```\n\nExplanation:\n\n``` \"A\" wins, he always plays first.\n\"X \" \"X \" \"X \" \"X \" \"```\n\nX\n\n``` \"\n\" \" -> \" \" -> \" X \" -> \" X \" -> \"```\n\nX\n\n``` \"\n\" \" \"O \" \"O \" \"OO \" \"OO```\n\nX\n\n```\"\n```\n\nExample 2:\n\nInput:\n\n``` moves = [[0,0],[1,1],[0,1],[0,2],[1,0],[2,0]]\n```\n\nOutput:\n\n``` \"B\"\n```\n\nExplanation:\n\n``` \"B\" wins.\n\"X \" \"X \" \"XX \" \"XXO\" \"XXO\" \"XX```\n\nO\n\n```\"\n\" \" -> \" O \" -> \" O \" -> \" O \" -> \"XO \" -> \"X```\n\nO\n\n``` \"\n\" \" \" \" \" \" \" \" \" \" \"```\n\nO\n\n``` \"\n```\n\nExample 3:\n\nInput:\n\n``` moves = [[0,0],[1,1],[2,0],[1,0],[1,2],[2,1],[0,1],[0,2],[2,2]]\n```\n\nOutput:\n\n``` \"Draw\"\n```\n\nExplanation:\n\n``` The game ends in a draw since there are no moves to make.\n\"XXO\"\n\"OOX\"\n\"XOX\"\n```\n\nExample 4:\n\nInput:\n\n``` moves = [[0,0],[1,1]]\n```\n\nOutput:\n\n``` \"Pending\"\n```\n\nExplanation:\n\n``` The game has not finished yet.\n\"X \"\n\" O \"\n\" \"\n```\n\nConstraints:\n\n• `1 <= moves.length <= 9`\n• `moves[i].length == 2`\n• `0 <= moves[i][j] <= 2`\n• There are no repeated elements on `moves`.\n• `moves` follow the rules of tic tac toe.\n```class Solution {\npublic String tictactoe(int[][] moves) {\n\nfor (int a = 0; a <= 1; a++ ){\nint[] x = new int;\nint[] y = new int;\nint[] v = new int;\n\nfor (int i = a ; i < moves.length; i += 2){\nx[moves[i]] ++;\ny[moves[i]] ++;\nif (moves[i] == moves[i]){\nv ++;\n}\nif (moves[i] == 2 - moves[i]) {\nv ++;\n}\n}\n\nif (x == 3 || x == 3 ||x == 3 || y == 3 || y == 3|| y == 3 || v == 3 || v == 3){\nreturn a==0? \"A\":\"B\";\n}\n}\n\nreturn moves.length == 9? \"Draw\":\"Pending\";\n}\n}```\n\n# 1266. Minimum Time Visiting All Points\n\nOn a plane there are `n` points with integer coordinates `points[i] = [xi, yi]`. Your task is to find the minimum time in seconds to visit all points.\n\nYou can move according to the next rules:\n\n• In one second always you can either move vertically, horizontally by one unit or diagonally (it means to move one unit vertically and one unit horizontally in one second).\n• You have to visit the points in the same order as they appear in the array.\n\nExample 1:", null, "Input:\n\n``` points = [[1,1],[3,4],[-1,0]]\n```\n\nOutput:\n\n``` 7\n```\n\nExplanation:\n\n`One optimal path is`\n\n[1,1]\n\n` -> [2,2] -> [3,3] ->`\n\n[3,4]\n\n`-> [2,3] -> [1,2] -> [0,1] ->`\n\n[-1,0]\n\n```\nTime from [1,1] to [3,4] = 3 seconds\nTime from [3,4] to [-1,0] = 4 seconds\nTotal time = 7 seconds```\n\nExample 2:\n\nInput:\n\n``` points = [[3,2],[-2,2]]\n```\n\nOutput:\n\n``` 5\n```\n\nConstraints:\n\n• `points.length == n`\n• `1 <= n <= 100`\n• `points[i].length == 2`\n• `-1000 <= points[i], points[i] <= 1000`\n\n```class Solution {\npublic int minTimeToVisitAllPoints(int[][] points) {\nint result = 0;\nint prev_x = points;\nint prev_y = points;\nfor (int i = 1; i < points.length; i ++){\nint x = points[i];\nint y = points[i];\n\nint diff_x = Math.abs(x - prev_x);\nint diff_y = Math.abs(y - prev_y);\n\nresult += Math.min(diff_x,diff_y)+Math.abs(diff_x-diff_y);\n\nprev_x = x;\nprev_y = y;\n}\nreturn result;\n}\n}```\n\n# 189. Rotate Array\n\nGiven an array, rotate the array to the right by k steps, where k is non-negative.\n\nExample 1:\n\nInput:\n\n``[1,2,3,4,5,6,7]` and`\n\nk\n\n` = 3`\n\nOutput:\n\n``[5,6,7,1,2,3,4]``\n\nExplanation:\n\n```rotate 1 steps to the right: `[7,1,2,3,4,5,6]` rotate 2 steps to the right: ```[6,7,1,2,3,4,5]\n```rotate 3 steps to the right: `[5,6,7,1,2,3,4]````\n\nExample 2:\n\nInput:\n\n``[-1,-100,3,99]` and`\n\nk\n\n` = 2`\n\nOutput:\n\n` [3,99,-1,-100]`\n\nExplanation:\n\n` rotate 1 steps to the right: [99,-1,-100,3] rotate 2 steps to the right: [3,99,-1,-100]`\n\nNote:\n\n• Try to come up as many solutions as you can, there are at least 3 different ways to solve this problem.\n• Could you do it in-place with O(1) extra space?\n\nIn place rotate\n\n```class Solution {\n\n/**\n* @param Integer[] \\$nums\n* @param Integer \\$k\n* @return NULL\n*/\nfunction rotate(&\\$nums, \\$k) {\n\\$c = count(\\$nums);\n\\$steps = 0;\n\\$start = 0;\n\nwhile(\\$steps < \\$c){\n\\$i = \\$start;\n\\$prev = \\$nums[\\$i];\ndo{\n\\$next = (\\$i+\\$k)%\\$c;\n\\$temp = \\$nums[\\$next];\n\\$nums[\\$next] = \\$prev;\n\\$prev = \\$temp;\n\\$i += \\$k;\n\\$i = \\$i % \\$c;\n\\$steps++;\n}while(\\$start != \\$i);\n\\$start ++;\n\n}\n\n}\n}```\n\n# 917. Reverse Only Letters\n\nGiven a string `S`, return the “reversed” string where all characters that are not a letter stay in the same place, and all letters reverse their positions.\n\nExample 1:\n\nInput:\n\n```\"ab-cd\"\n```\n\nOutput:\n\n```\"dc-ba\"\n```\n\nExample 2:\n\nInput:\n\n```\"a-bC-dEf-ghIj\"\n```\n\nOutput:\n\n```\"j-Ih-gfE-dCba\"\n```\n\nExample 3:\n\nInput:\n\n```\"Test1ng-Leet=code-Q!\"\n```\n\nOutput:\n\n```\"Qedo1ct-eeLg=ntse-T!\"\n```\n\nNote:\n\n1. `S.length <= 100`\n2. `33 <= S[i].ASCIIcode <= 122`\n3. `S` doesn’t contain `\\` or `\"`\n\nStraightforward solution.\n\n1. collect letters in reverse order.\n2. go through string from beginning, if it’s letter, replace with letter from reverse string, else keep special chars and adjust pointer accordingly\n3. of course there’s space for improvement, but not necessary\n```<?php\nclass Solution {\n\n/**\n* @param String \\$S\n* @return String\n*/\nfunction reverseOnlyLetters(\\$S) {\n\\$x = '';\n\\$y = '';\nfor(\\$i = strlen(\\$S)-1; \\$i >=0; \\$i -- ){\nif(ctype_alpha(\\$S[\\$i])){\n\\$x .= \\$S[\\$i];\n}\n}\n\\$xi = 0;\nfor(\\$i = 0; \\$i < strlen(\\$S); \\$i ++){\nif(!ctype_alpha(\\$S[\\$i])){\n\\$y .= \\$S[\\$i];\n} else {\n\\$y .= \\$x[\\$xi];\n\\$xi ++;\n}\n\n}\n\nreturn \\$y;\n}\n}\n\n?>```\n\n# 7. Reverse Integer\n\nGiven a 32-bit signed integer, reverse digits of an integer.\n\nExample 1:\n\nInput:\n\n``` 123\n```\n\nOutput:\n\n``` 321\n```\n\nExample 2:\n\nInput:\n\n``` -123\n```\n\nOutput:\n\n``` -321\n```\n\nExample 3:\n\nInput:\n\n``` 120\n```\n\nOutput:\n\n``` 21\n```\n\nNote:\nAssume we are dealing with an environment which could only store integers within the 32-bit signed integer range: [−231,  231 − 1]. For the purpose of this problem, assume that your function returns 0 when the reversed integer overflows.\n\n```class Solution {\n\n/**\n* @param Integer \\$x\n* @return Integer\n*/\nfunction reverse(\\$x) {\n\n\\$r = (int) (strrev(trim(\\$x,\"+-\")) * (\\$x>0?1:-1));\n\nif(\\$r > pow(2,31)-1 || \\$r < -pow(2,31)){\nreturn 0;\n}\nreturn \\$r;\n}\n}```" ]
[ null, "https://assets.leetcode.com/uploads/2019/11/14/1626_example_1.PNG", null ]
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https://chemistry.stackexchange.com/questions/81862/derivation-of-formula-used-in-calculating-molecular-mass-from-a-mixture
[ "# Derivation of formula used in calculating molecular mass from a mixture\n\nIn a solution to a problem, to calculate Molecular mass of unknown gas from a mixture containing $\\ce{O2}$ in $80~\\%$, the author used the formula :\n\n$\\frac{1}{\\sqrt{M_\\text{mix}}}=\\frac{X_{\\ce{O2}}}{\\sqrt{M_{\\ce{O2}}}}+\\frac{X_\\text{gas}}{\\sqrt{M_\\text{gas}}}$.\n\nThe exact problem is:\n\nPure $\\ce{O_2}$ diffuses through an aperture in $224\\mathrm{s}$, where as mixture of $\\ce{O_2}$ and another gas containing $80~\\%\\ \\ce{O_2}$ diffuses from the same in $234\\mathrm{s}$. The molecular mass of gas will be?\n\nSolution by author:\n\n\\begin{align}\\frac{t_\\text{mix}}{t_{\\ce{O2}}} &= \\frac{r_{\\ce{O2}}}{r_\\text{mix}} = \\sqrt{\\frac{M_\\text{mix}}{32}}\\\\[1em] \\frac{234}{224} &= \\sqrt{\\frac{M_\\text{mix}}{32}}\\\\[0.5em] M_\\text{mix} &= 34.92\\\\[1em] \\Longrightarrow \\qquad \\frac{1}{\\sqrt{M_\\text{mix}}} &= \\frac{X_\\text{gas}}{\\sqrt{M_\\text{gas}}} + \\frac{X_{\\ce{O2}}}{\\sqrt{M_{\\ce{O2}}}}\\\\[1em] \\Longrightarrow \\qquad \\frac{1}{\\sqrt{34.92}} &= \\frac{0.2}{\\sqrt{M_\\text{gas}}} + \\frac{0.8}{\\sqrt{32}}\\\\[0.5em] M_\\text{gas} &= 51.5\\end{align}\n\nI wish to know why this works here?\n\nHint: $$v_{rms}=\\sqrt{\\frac{3RT}M}$$\n\nAnd weighted average of speeds for mixture.\n\nEDIT:\n\n$$v_{mix} ^{rms}= x_1 v_1^{rms}+ x_2 v_2^{rms}$$\n\nThis will complete the derivation of the formula you gave.\n\n• This does not provide an answer to the question. To critique or request clarification from an author, leave a comment below their post. - From Review – andselisk Aug 29 '17 at 6:07\n• Actually there is just one step to the derivation after you write down both the equations.. I guess OP can solve that on his own. I think spoon feeding complete derivation won't do any good. – ABC Aug 29 '17 at 6:30\n• Hint answers aren't really well-received here. If you want to give a hint, it's better to do it as a comment. – orthocresol Aug 29 '17 at 6:31" ]
[ null ]
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https://melindadaily.com/skin-problem/how-many-electrons-are-in-a-mole-of-aluminum.html
[ "# How many electrons are in a mole of aluminum?\n\nContents\n\nAluminum has 13 electrons. The atomic number of an element tells us how many protons and electrons are in an atom of that element. Since the atomic…\n\n## How many electrons are in a mole?\n\n1 mole of electrons contains the Avogadro constant, L, electrons – that is 6.02 x 1023 electrons. You would also be given that in an exam if you needed to use it. This value is known as the Faraday constant.\n\n## Can aluminum have 12 electrons?\n\nAluminum has 13 electrons in a neutral atom. … The 1S2, 2S2 and 2P6 shells are full, and so is the 3S2. The 3P1 shell has one electron in it (in Al’s neutral atom) and that gives aluminum one lone valence electron and two “semi-valence” electrons which can be loaned out.\n\n## How do you find the number of electrons in mole concept?\n\nExplanation: First, we calculate the number of electrons in ONE MOLECULE of CO2 . There is ONE CARBON ATOM, that is 6 electrons; and TWO OXYGEN ATOMS, that is 16 electrons, i.e. 22 electrons per molecule. Moles of carbon dioxide = 100⋅g44.01⋅g⋅mol−1=2.27⋅mol .\n\n## Why does aluminum have 13 electrons?\n\nThe electrons like to be in separate shells/orbitals. Shell number one can only hold 2 electrons, shell two can hold 8, and for the first eighteen elements shell three can hold a maximum of eight electrons. … That means there are 13 electrons in a aluminum atom.\n\n## How many electrons are in a mole of Sc3+?\n\nOxygen has (on average) 8 protons, neutrons and electrons, for a total of 24, and H20 has a total number of 28 particles. Therfore there are 18 electrons,20 neutrons and 20 protons present in calcium ion. 7 electrons, 7 protons, and 7 neutrons. The Sc3+ ion contains 21 protons and 18 electrons.\n\n## Can you have a mole of electrons?\n\nSince 1 mol electrons = 6.02214076×1023 electrons (Avogadro’s number), the Faraday constant is equal to the elementary charge e, the magnitude of the charge of an electron, multiplied by 1 mole: F = 96485.\n\n## Why does aluminum have a charge of 3+?\n\nWhat is the charge of aluminum ion? … As Al has 3 valence electrons as M(3) so it tends to lose electrons then after losing these 3 electrons it has 10 electrons and 13 protons, and 10 electrons are neutralized by 10 protons out of 13, and the further excess 3 protons out of 13 appear as +3 charge.\n\n## Why is Aluminium electrically neutral?\n\nPositively-charged protons attract negatively-charged electrons, but the negatively-charged electrons repel one-another. As a result, the atom can attract a number of electrons until it has equal numbers of protons and electrons, making the atom neutral.\n\n13\n\n## How do you calculate the number of electrons?\n\nThe number of protons in the nucleus of the atom is equal to the atomic number (Z). The number of electrons in a neutral atom is equal to the number of protons. The mass number of the atom (M) is equal to the sum of the number of protons and neutrons in the nucleus.\n\nIT IS INTERESTING:  Is Neutrogena Hydro Boost Gel Cream good for acne prone skin?\n\n## How many electrons are in a gram?\n\nHow many electrons in 1 gram? The answer is 1.0977693108051E+27. We assume you are converting between electron and gram.\n\n## How many valence electrons are in 1.523 moles of CH4?\n\nValence Electrons And 8 One CHa Molecule Has 8 Valence Electrons Are In 1.523 Moles Of CH4 (Answers In The Form Of 1.234E12: Avogadro’s Number 6.022E23) Total Electrons Are In 1.523 Moles Of CH4 MM: 16,04 G/mol 1.523 Moles Of CHe Has A Mass Of 16.04246 Click Save And Submit To Save And Submit. …\n\n## Can aluminum be picked up by a magnet?\n\nIn our everyday experience aluminum doesn’t stick to magnets (neither does copper). Most matter will exhibit some magnetic attraction when under high enough magnetic fields. … But under normal circumstances aluminum isn’t visibly magnetic.\n\n## What does aluminum need to become stable?\n\nAluminum is a metal that will always lose three electrons. The halogens all have seven valence electrons. Each one of these elements wants to gain one electron to achieve an octet.\n\n## Why is aluminum so easily recycled?\n\nRecycling aluminum saves more than 90 percent of the energy needed to make new aluminum. Recycling aluminum saves more than 90 percent of the energy that would be needed to create a comparable amount of the metal from raw materials.", null, "" ]
[ null, "data:image/svg+xml,%3Csvg%20xmlns='http://www.w3.org/2000/svg'%20viewBox='0%200%201000%20947'%3E%3C/svg%3E", null ]
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https://www.colorhexa.com/02ba2e
[ "# #02ba2e Color Information\n\nIn a RGB color space, hex #02ba2e is composed of 0.8% red, 72.9% green and 18% blue. Whereas in a CMYK color space, it is composed of 98.9% cyan, 0% magenta, 75.3% yellow and 27.1% black. It has a hue angle of 134.3 degrees, a saturation of 97.9% and a lightness of 36.9%. #02ba2e color hex could be obtained by blending #04ff5c with #007500. Closest websafe color is: #00cc33.\n\n• R 1\n• G 73\n• B 18\nRGB color chart\n• C 99\n• M 0\n• Y 75\n• K 27\nCMYK color chart\n\n#02ba2e color description : Strong lime green.\n\n# #02ba2e Color Conversion\n\nThe hexadecimal color #02ba2e has RGB values of R:2, G:186, B:46 and CMYK values of C:0.99, M:0, Y:0.75, K:0.27. Its decimal value is 178734.\n\nHex triplet RGB Decimal 02ba2e `#02ba2e` 2, 186, 46 `rgb(2,186,46)` 0.8, 72.9, 18 `rgb(0.8%,72.9%,18%)` 99, 0, 75, 27 134.3°, 97.9, 36.9 `hsl(134.3,97.9%,36.9%)` 134.3°, 98.9, 72.9 00cc33 `#00cc33`\nCIE-LAB 66.002, -65.92, 56.071 18.076, 35.326, 8.45 0.292, 0.571, 35.326 66.002, 86.542, 139.616 66.002, -61.541, 73.974 59.436, -49.726, 33.175 00000010, 10111010, 00101110\n\n# Color Schemes with #02ba2e\n\n• #02ba2e\n``#02ba2e` `rgb(2,186,46)``\n• #ba028e\n``#ba028e` `rgb(186,2,142)``\nComplementary Color\n• #32ba02\n``#32ba02` `rgb(50,186,2)``\n• #02ba2e\n``#02ba2e` `rgb(2,186,46)``\n• #02ba8a\n``#02ba8a` `rgb(2,186,138)``\nAnalogous Color\n• #ba0232\n``#ba0232` `rgb(186,2,50)``\n• #02ba2e\n``#02ba2e` `rgb(2,186,46)``\n• #8a02ba\n``#8a02ba` `rgb(138,2,186)``\nSplit Complementary Color\n• #ba2e02\n``#ba2e02` `rgb(186,46,2)``\n• #02ba2e\n``#02ba2e` `rgb(2,186,46)``\n• #2e02ba\n``#2e02ba` `rgb(46,2,186)``\n• #8eba02\n``#8eba02` `rgb(142,186,2)``\n• #02ba2e\n``#02ba2e` `rgb(2,186,46)``\n• #2e02ba\n``#2e02ba` `rgb(46,2,186)``\n• #ba028e\n``#ba028e` `rgb(186,2,142)``\n• #016e1b\n``#016e1b` `rgb(1,110,27)``\n• #018822\n``#018822` `rgb(1,136,34)``\n• #02a128\n``#02a128` `rgb(2,161,40)``\n• #02ba2e\n``#02ba2e` `rgb(2,186,46)``\n• #02d334\n``#02d334` `rgb(2,211,52)``\n• #03ec3a\n``#03ec3a` `rgb(3,236,58)``\n• #0cfc46\n``#0cfc46` `rgb(12,252,70)``\nMonochromatic Color\n\n# Alternatives to #02ba2e\n\nBelow, you can see some colors close to #02ba2e. Having a set of related colors can be useful if you need an inspirational alternative to your original color choice.\n\n• #04ba02\n``#04ba02` `rgb(4,186,2)``\n• #02ba0f\n``#02ba0f` `rgb(2,186,15)``\n• #02ba1f\n``#02ba1f` `rgb(2,186,31)``\n• #02ba2e\n``#02ba2e` `rgb(2,186,46)``\n• #02ba3d\n``#02ba3d` `rgb(2,186,61)``\n• #02ba4d\n``#02ba4d` `rgb(2,186,77)``\n• #02ba5c\n``#02ba5c` `rgb(2,186,92)``\nSimilar Colors\n\n# #02ba2e Preview\n\nThis text has a font color of #02ba2e.\n\n``<span style=\"color:#02ba2e;\">Text here</span>``\n#02ba2e background color\n\nThis paragraph has a background color of #02ba2e.\n\n``<p style=\"background-color:#02ba2e;\">Content here</p>``\n#02ba2e border color\n\nThis element has a border color of #02ba2e.\n\n``<div style=\"border:1px solid #02ba2e;\">Content here</div>``\nCSS codes\n``.text {color:#02ba2e;}``\n``.background {background-color:#02ba2e;}``\n``.border {border:1px solid #02ba2e;}``\n\n# Shades and Tints of #02ba2e\n\nA shade is achieved by adding black to any pure hue, while a tint is created by mixing white to any pure color. In this example, #000b03 is the darkest color, while #f7fff9 is the lightest one.\n\n• #000b03\n``#000b03` `rgb(0,11,3)``\n• #001f08\n``#001f08` `rgb(0,31,8)``\n• #01320c\n``#01320c` `rgb(1,50,12)``\n• #014611\n``#014611` `rgb(1,70,17)``\n• #015916\n``#015916` `rgb(1,89,22)``\n• #016c1b\n``#016c1b` `rgb(1,108,27)``\n• #018020\n``#018020` `rgb(1,128,32)``\n• #029324\n``#029324` `rgb(2,147,36)``\n• #02a729\n``#02a729` `rgb(2,167,41)``\n• #02ba2e\n``#02ba2e` `rgb(2,186,46)``\n• #02cd33\n``#02cd33` `rgb(2,205,51)``\n• #02e138\n``#02e138` `rgb(2,225,56)``\n• #03f43c\n``#03f43c` `rgb(3,244,60)``\n• #0efc47\n``#0efc47` `rgb(14,252,71)``\n• #21fd56\n``#21fd56` `rgb(33,253,86)``\n• #35fd65\n``#35fd65` `rgb(53,253,101)``\n• #48fd73\n``#48fd73` `rgb(72,253,115)``\n• #5cfd82\n``#5cfd82` `rgb(92,253,130)``\n• #6ffd91\n``#6ffd91` `rgb(111,253,145)``\n• #82fea0\n``#82fea0` `rgb(130,254,160)``\n• #96feaf\n``#96feaf` `rgb(150,254,175)``\n• #a9febe\n``#a9febe` `rgb(169,254,190)``\n• #bdfecc\n``#bdfecc` `rgb(189,254,204)``\n• #d0fedb\n``#d0fedb` `rgb(208,254,219)``\n• #e4ffea\n``#e4ffea` `rgb(228,255,234)``\n• #f7fff9\n``#f7fff9` `rgb(247,255,249)``\nTint Color Variation\n\n# Tones of #02ba2e\n\nA tone is produced by adding gray to any pure hue. In this case, #59635b is the less saturated color, while #02ba2e is the most saturated one.\n\n• #59635b\n``#59635b` `rgb(89,99,91)``\n• #526a57\n``#526a57` `rgb(82,106,87)``\n• #4a7254\n``#4a7254` `rgb(74,114,84)``\n• #437950\n``#437950` `rgb(67,121,80)``\n• #3c804c\n``#3c804c` `rgb(60,128,76)``\n• #358748\n``#358748` `rgb(53,135,72)``\n• #2d8f45\n``#2d8f45` `rgb(45,143,69)``\n• #269641\n``#269641` `rgb(38,150,65)``\n• #1f9d3d\n``#1f9d3d` `rgb(31,157,61)``\n• #18a439\n``#18a439` `rgb(24,164,57)``\n• #10ac36\n``#10ac36` `rgb(16,172,54)``\n• #09b332\n``#09b332` `rgb(9,179,50)``\n• #02ba2e\n``#02ba2e` `rgb(2,186,46)``\nTone Color Variation\n\n# Color Blindness Simulator\n\nBelow, you can see how #02ba2e is perceived by people affected by a color vision deficiency. This can be useful if you need to ensure your color combinations are accessible to color-blind users.\n\nMonochromacy\n• Achromatopsia 0.005% of the population\n• Atypical Achromatopsia 0.001% of the population\nDichromacy\n• Protanopia 1% of men\n• Deuteranopia 1% of men\n• Tritanopia 0.001% of the population\nTrichromacy\n• Protanomaly 1% of men, 0.01% of women\n• Deuteranomaly 6% of men, 0.4% of women\n• Tritanomaly 0.01% of the population" ]
[ null ]
{"ft_lang_label":"__label__en","ft_lang_prob":0.5293136,"math_prob":0.77119327,"size":3666,"snap":"2019-51-2020-05","text_gpt3_token_len":1634,"char_repetition_ratio":0.12834516,"word_repetition_ratio":0.011090573,"special_character_ratio":0.54964536,"punctuation_ratio":0.23608018,"nsfw_num_words":0,"has_unicode_error":false,"math_prob_llama3":0.9874096,"pos_list":[0],"im_url_duplicate_count":[null],"WARC_HEADER":"{\"WARC-Type\":\"response\",\"WARC-Date\":\"2019-12-16T01:00:55Z\",\"WARC-Record-ID\":\"<urn:uuid:f4b81187-5923-4466-a35c-dc9b69b70944>\",\"Content-Length\":\"36218\",\"Content-Type\":\"application/http; msgtype=response\",\"WARC-Warcinfo-ID\":\"<urn:uuid:7c3bf365-c9a4-46c9-87c6-48dc70f9dcd4>\",\"WARC-Concurrent-To\":\"<urn:uuid:8fbb1330-729b-4253-a6f9-f3e4f6b280d1>\",\"WARC-IP-Address\":\"178.32.117.56\",\"WARC-Target-URI\":\"https://www.colorhexa.com/02ba2e\",\"WARC-Payload-Digest\":\"sha1:NKUTEAXQVX5CMCOB5TWN3NWEYXNTIGUP\",\"WARC-Block-Digest\":\"sha1:7CFDGSAIMATZP4QNUWBKHFJYTGCLLAQG\",\"WARC-Identified-Payload-Type\":\"text/html\",\"warc_filename\":\"/cc_download/warc_2019/CC-MAIN-2019-51/CC-MAIN-2019-51_segments_1575541310970.85_warc_CC-MAIN-20191215225643-20191216013643-00158.warc.gz\"}"}
https://blog.newtum.com/addition-of-two-numbers-in-c/
[ "# Addition of Two Numbers in C\n\n(Last Updated On: 13/09/2023)\n\nThis program explains how to write the addition of two numbers in c. For an explanation, we have used only integer variables. This program accepts two inputs using the scanf input function. Then adds the two numbers and store them in the third variable. The output is displayed using the printf statement.\n\n```//SUM of two numbers in C Programming\nmain()\n{\nint a,b,c;\nprintf(\"\\nEnter value for a : \");\nscanf(\"%d\", &a);\nprintf(\"\\nEnter value for b : \");\nscanf(\"%d\", &b);\nc=a+b;\nprintf(\"Result : %d\",c);\n}```\n\nAdding two numbers in C is the primary program taught while learning C Programming in each and every C Programming course. Adding two numbers in C teacher’s basics of variable declaration and very basic of mathematical operators.\n\nC Program to Add Two Numbers explains the use of basic functions of input and output (scanf and printf) in programming. This blog demonstrates two methods to add two numbers in C programming.\n\nOne is a static method, and the other is dynamic, where the program accepts input from the user. To learn more about C Programming and its history, please read this blog first, “C Programming Complete Guide”. You can watch the video over here, which explains the blog in a simple video with VFX animation.\n\nIf you want to see the complete course with all the topics of C programming covered, you can refer to this link for C Programming with Certification.\n\n## C Program to Add Two Numbers using scanf() Statement\n\nHow to achieve this. Now we will be learning a new function called scanf(). scanf() takes the input from the users and does the processing.\n\nDeclaration of Variables\n\nHere we will declare the a,b, and c variables in one line. Since no initialization is required, we can declare the variables in one line. Since we are assuming all the values to be integers, we will use the word int.\n\nAccepting User Inputs using scanf statement\n\nC uses a scanf statement to accept two numbers. Since we are accepting input in variables a,b and both are integers, our input statement to add two numbers in C Programming will look like the below code.\n\n#### Simple Logic to add two numbers in C Programming\n\nSimple logic is very, very easy to add two numbers. First, We will write a small mathematical expression “c = a+ b”, and then we will use the printf statement to display the final result.\n\n```main()\n{\nint a,b,c;\nscanf(\"%d\", &a);\nscanf(\"%d\", &b);\nc=a+b;\nprint(\"%d\",c);\n}```\n##### Output\n\nLet’s look at the output. When we execute the program, a black output screen will appear. It’s because of a scanf. It’s waiting for our input. Now since we have specified scanf twice, it will wait for two inputs. Let’s give 21 as input. One input is gone, and now the program is waiting for the second input. Now enter value 4 and again press enter. You will see 25 as output.\n\nOnly scanf waits for input. Once the input is given, the program runs the expression a+b, stores the result in c, and then prints the value of c.\n\nBefore understanding the scanf and this program in detail, let’s modify it so that it’s more user-friendly. After declaring variables, we write a printf statement to give a title and another line to ask the user to give the first input. Same for the second input.\n\n```main()\n{\nint a,b,c;\nprintf(\"\\nEnter value for a : \");\nscanf(\"%d\", &a);\nprintf(\"\\nEnter value for b : \");\nscanf(\"%d\", &b);\nc=a+b;\nprintf(\"Result : %d\",c);\n}```\n\nLet’s execute the program. As we execute it, the system displays the title first. You can view the source code from the Newtum GitHub repository.\n\n##### Output :\n``````Welcome to program of Addition\nEnter value for a : 12\nEnter value for b : 18\nResult : 30``````\n\nI hope you got some more understanding of the printf and business logic or problem definition in the C programming statement. Let’s explore the scanf statement. I am saying again that scanf takes the input. It has two parameters, the First is placeholders of variables, and the second is the variable followed by an ampersand (&). Ampersand is the biggest and most difficult part of the C language.\n\nFor now, just assume that it’s a syntax, I mean a way for writing a scanf statement. We have used the printf statement to give the output and the scanf statement to take input. You can download the code from here.\n\nFor Other Countries: To watch C Programming Content in another country, please visit the URL:https://newtum.com/course-details/c-programming-online.\n\n## Add Two Numbers in C without using scanf statement\n\nThis method doesn’t use the scanf function and embeds hardcoded values in the variable. To change the program, we have to edit the program to put new values. Though this method is easy, it’s not good for professional C programming. Since we need to understand this concept, we have explained this method.\n\nDeclaring Variables in C Program\n\nHere, to start, we will create a program that will add two numbers. Here we will create 3 integers using the keyword int variables a, b, and c. We want to tell you that we can create as many variables in a single line while using int keywords only one time.\n\nInitialize variable\n\nNow we put value 10 in a and value 20 in b. Initialization is a process in C Programming where we initialize variables at the start of the program or business logic. Though we have used it to add two numbers in C Programming, this is mainly used in loops in C.\n\n#### Simple Logic to add two numbers in C Programming\n\nBusiness logic is very, very easy to add two numbers. First, We will write a small mathematical expression “c = a+ b”, and then we will use the printf statement to display the final result.\n\n```main()\n{\nint a=10;\nint b=20;\nint c=a+b;\nprintf(\"%d\",c);\n}```\n\nPlease make sure that all the statements are followed by a semicolon. I have explained the reason in the very first chapter. Now we print the variable c. Execute the program.\n\n##### Output :\n``30``\n\nSo what have we done? We have created 3 variables. A and b act as input variables, while C stores the result of an expression, i.e. a+b, and we print the result. It’s a pretty simple program. But we can’t have fixed values, i.e. a=10 and b=10. And everyone doesn’t know programming languages. We need to prepare a program that can be used by anyone. We don’t have to write values of a and b via editing the program.\n\nWhat is the important function required to write a program to add two numbers in C\n\nFor the addition of two numbers, the scanf and printf function is important.\n\n###### scanf():\n• scanf() function is used to take input from the user.\n• This function allows you to accept input from standard in, which for us is generally the keyboard.\n• The simplest application of scanf looks like this: scanf(“%d”, &a);\n###### print():\n• printf() function is used to display output to the user.\n• printf() function is used to print the “character, string, float, integer values” onto the output screen.\n\nNumber of a variable required to add two numbers in C Programming\n\nAt least three variables are required for the addition of two numbers. In which two variables are used for the stored integer value and a third variable for the result." ]
[ null ]
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https://erikbern.com/2014/12/11/deep-learning-for-go.html
[ "# Deep learning for&#8230; Go\n\nThis is the last post about deep learning for chess/go/whatever. But this really cool paper by Christopher Clark and Amos Storkey was forwarded to me by Michael Eickenberg. It's about using convolutional neural networks to play Go. The authors of the paper do a much better job than I would ever have done of modeling move prediction in Go and show that their model beat certain Go engines.\n\nThe fascinating thing about this paper is that playing against other Go engines, they just plug in their move prediction function, with no deep search beyond one level. That means the total time it spends is a fraction of the opponents. Still, the fact that it plays so well speaks for its strength.\n\nSo what happened if we could plug this into a deep search framework? The authors suggest doing exactly that in the conclusion. State of the art of Go engines actually use Monte Carlo tree search rather than minimax but other than that, it's the same principle.\n\nI talked a bit with the authors and the main thing that you have to change is to switch from move prediction to an evaluation function. For my chess experiments, I found a (hacky) way to train a function that does both at the same time. There's essentially two terms in my objective function: one is comparing the actual move with a random move, using a sigmoid:\n\n$$frac{P(q)}{P(q) + P(r)} = frac{exp(f(q))}{exp(f(q)) + exp(f(r))}$$ .\n\nIf you extend that to all possible random moves you actually get a full probability distribution (a softmax) over all possible next moves.\n\n$$P(p rightarrow q) = frac{exp(f(q))}{sum exp(f(r)) }$$ .\n\nNow, how do you “convert” that into an evaluation function? That's the second term, which tries fit the negative parent score to the current score. We penalize the quantity $$f(p) + f(q)$$ by throwing in two more sigmoids. It's a “soft constraint” that has absolutely no probabilistic interpretation. This a hacky solution, but here's how I justify it:\n\n1. Note that the evaluation functions are unique up to a monotonic transform, so we can actually mangle it quite a lot.\n2. The softmax distribution has one degree of freedom in how it chooses the quantities, so (I'm speculating) the artificial constraint does not change the probabilities.\n\nI think you could do the exact thing with their Go engine. In fact I'm willing to bet a couple of hundred bucks that if you did that, you would end up with the best Go engine in the world.\n\nBtw another fun thing was that they plot some of the filters and they seem as random as the ones I learn for Chess. But a clever trick enforcing symmetry seem to help the model quite a lot.", null, "Tagged with: math" ]
[ null, "https://erikbern.com/assets/2014/12/Screen-Shot-2014-12-11-at-5.11.32-PM.png", null ]
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https://kornia.readthedocs.io/en/latest/losses.html
[ "# kornia.losses¶\n\n## Reconstruction¶\n\nssim(img1: torch.Tensor, img2: torch.Tensor, window_size: int, reduction: str = 'none', max_val: float = 1.0) → torch.Tensor[source]\n\nFunction that measures the Structural Similarity (SSIM) index between each element in the input x and target y.\n\nSee SSIM for details.\n\npsnr_loss(input: torch.Tensor, target: torch.Tensor, max_val: float) → torch.Tensor[source]\n\nFunction that computes PSNR\n\nSee PSNRLoss for details.\n\ntotal_variation(img: torch.Tensor) → torch.Tensor[source]\n\nFunction that computes Total Variation.\n\nSee TotalVariation for details.\n\ninverse_depth_smoothness_loss(idepth: torch.Tensor, image: torch.Tensor) → torch.Tensor[source]\n\nComputes image-aware inverse depth smoothness loss.\n\nSee InvDepthSmoothnessLoss for details.\n\n## Semantic Segmentation¶\n\nfocal_loss(input: torch.Tensor, target: torch.Tensor, alpha: float, gamma: float = 2.0, reduction: str = 'none', eps: float = 1e-08) → torch.Tensor[source]\n\nFunction that computes Focal loss.\n\nSee FocalLoss for details.\n\ndice_loss(input: torch.Tensor, target: torch.Tensor, eps: float = 1e-08) → torch.Tensor[source]\n\nFunction that computes Sørensen-Dice Coefficient loss.\n\nSee DiceLoss for details.\n\ntversky_loss(input: torch.Tensor, target: torch.Tensor, alpha: float, beta: float, eps: float = 1e-08) → torch.Tensor[source]\n\nFunction that computes Tversky loss.\n\nSee TverskyLoss for details.\n\n## Distributions¶\n\njs_div_loss_2d(input: torch.Tensor, target: torch.Tensor, reduction: str = 'mean')[source]\n\nCalculates the Jensen-Shannon divergence loss between heatmaps.\n\nParameters\n• input (torch.Tensor) – the input tensor.\n\n• target (torch.Tensor) – the target tensor.\n\n• reduction (string, optional) – Specifies the reduction to apply to the output: 'none' | 'mean' | 'sum'. 'none': no reduction will be applied, 'mean': the sum of the output will be divided by the number of elements in the output, 'sum': the output will be summed. Default: 'mean'.\n\nShape:\n• Input: $$(B, N, H, W)$$\n\n• Target: $$(B, N, H, W)$$, same shape as the input\n\nkl_div_loss_2d(input: torch.Tensor, target: torch.Tensor, reduction: str = 'mean')[source]\n\nCalculates the Kullback-Leibler divergence loss between heatmaps.\n\nParameters\n• input (torch.Tensor) – the input tensor.\n\n• target (torch.Tensor) – the target tensor.\n\n• reduction (string, optional) – Specifies the reduction to apply to the output: 'none' | 'mean' | 'sum'. 'none': no reduction will be applied, 'mean': the sum of the output will be divided by the number of elements in the output, 'sum': the output will be summed. Default: 'mean'.\n\nShape:\n• Input: $$(B, N, H, W)$$\n\n• Target: $$(B, N, H, W)$$, same shape as the input\n\n## Module¶\n\nclass DiceLoss[source]\n\nCriterion that computes Sørensen-Dice Coefficient loss.\n\nAccording to , we compute the Sørensen-Dice Coefficient as follows:\n\n$\\text{Dice}(x, class) = \\frac{2 |X| \\cap |Y|}{|X| + |Y|}$\nwhere:\n• $$X$$ expects to be the scores of each class.\n\n• $$Y$$ expects to be the one-hot tensor with the class labels.\n\nthe loss, is finally computed as:\n\n$\\text{loss}(x, class) = 1 - \\text{Dice}(x, class)$\nShape:\n• Input: $$(N, C, H, W)$$ where C = number of classes.\n\n• Target: $$(N, H, W)$$ where each value is $$0 ≤ targets[i] ≤ C−1$$.\n\nExamples\n\n>>> N = 5 # num_classes\n>>> loss = kornia.losses.DiceLoss()\n>>> input = torch.randn(1, N, 3, 5, requires_grad=True)\n>>> target = torch.empty(1, 3, 5, dtype=torch.long).random_(N)\n>>> output = loss(input, target)\n>>> output.backward()\n\nclass TverskyLoss(alpha: float, beta: float, eps: float = 1e-08)[source]\n\nCriterion that computes Tversky Coeficient loss.\n\nAccording to , we compute the Tversky Coefficient as follows:\n\n$\\text{S}(P, G, \\alpha; \\beta) = \\frac{|PG|}{|PG| + \\alpha |P \\setminus G| + \\beta |G \\setminus P|}$\nwhere:\n• $$P$$ and $$G$$ are the predicted and ground truth binary labels.\n\n• $$\\alpha$$ and $$\\beta$$ control the magnitude of the penalties for FPs and FNs, respectively.\n\nNotes\n\n• $$\\alpha = \\beta = 0.5$$ => dice coeff\n\n• $$\\alpha = \\beta = 1$$ => tanimoto coeff\n\n• $$\\alpha + \\beta = 1$$ => F beta coeff\n\nShape:\n• Input: $$(N, C, H, W)$$ where C = number of classes.\n\n• Target: $$(N, H, W)$$ where each value is $$0 ≤ targets[i] ≤ C−1$$.\n\nExamples\n\n>>> N = 5 # num_classes\n>>> loss = kornia.losses.TverskyLoss(alpha=0.5, beta=0.5)\n>>> input = torch.randn(1, N, 3, 5, requires_grad=True)\n>>> target = torch.empty(1, 3, 5, dtype=torch.long).random_(N)\n>>> output = loss(input, target)\n>>> output.backward()\n\n\nReferences\n\nclass FocalLoss(alpha: float, gamma: float = 2.0, reduction: str = 'none')[source]\n\nCriterion that computes Focal loss.\n\nAccording to , the Focal loss is computed as follows:\n\n$\\text{FL}(p_t) = -\\alpha_t (1 - p_t)^{\\gamma} \\, \\text{log}(p_t)$\nwhere:\n• $$p_t$$ is the model’s estimated probability for each class.\n\nParameters\n• alpha (float) – Weighting factor $$\\alpha \\in [0, 1]$$.\n\n• gamma (float) – Focusing parameter $$\\gamma >= 0$$.\n\n• reduction (str, optional) – Specifies the reduction to apply to the output: ‘none’ | ‘mean’ | ‘sum’. ‘none’: no reduction will be applied, ‘mean’: the sum of the output will be divided by the number of elements in the output, ‘sum’: the output will be summed. Default: ‘none’.\n\nShape:\n• Input: $$(N, C, *)$$ where C = number of classes.\n\n• Target: $$(N, *)$$ where each value is $$0 ≤ targets[i] ≤ C−1$$.\n\nExamples\n\n>>> N = 5 # num_classes\n>>> kwargs = {\"alpha\": 0.5, \"gamma\": 2.0, \"reduction\": 'mean'}\n>>> loss = kornia.losses.FocalLoss(**kwargs)\n>>> input = torch.randn(1, N, 3, 5, requires_grad=True)\n>>> target = torch.empty(1, 3, 5, dtype=torch.long).random_(N)\n>>> output = loss(input, target)\n>>> output.backward()\n\n\nReferences\n\nclass SSIM(window_size: int, reduction: str = 'none', max_val: float = 1.0)[source]\n\nCreates a criterion that measures the Structural Similarity (SSIM) index between each element in the input x and target y.\n\nThe index can be described as:\n\n$\\text{SSIM}(x, y) = \\frac{(2\\mu_x\\mu_y+c_1)(2\\sigma_{xy}+c_2)} {(\\mu_x^2+\\mu_y^2+c_1)(\\sigma_x^2+\\sigma_y^2+c_2)}$\nwhere:\n• $$c_1=(k_1 L)^2$$ and $$c_2=(k_2 L)^2$$ are two variables to stabilize the division with weak denominator.\n\n• $$L$$ is the dynamic range of the pixel-values (typically this is $$2^{\\#\\text{bits per pixel}}-1$$).\n\nthe loss, or the Structural dissimilarity (DSSIM) can be finally described as:\n\n$\\text{loss}(x, y) = \\frac{1 - \\text{SSIM}(x, y)}{2}$\nParameters\n• window_size (int) – the size of the kernel.\n\n• max_val (float) – the dynamic range of the images. Default: 1.\n\n• reduction (str, optional) – Specifies the reduction to apply to the output: ‘none’ | ‘mean’ | ‘sum’. ‘none’: no reduction will be applied, ‘mean’: the sum of the output will be divided by the number of elements in the output, ‘sum’: the output will be summed. Default: ‘none’.\n\nReturns\n\nthe ssim index.\n\nReturn type\n\nTensor\n\nShape:\n• Input: $$(B, C, H, W)$$\n\n• Target $$(B, C, H, W)$$\n\n• Output: scale, if reduction is ‘none’, then $$(B, C, H, W)$$\n\nExamples:\n\n>>> input1 = torch.rand(1, 4, 5, 5)\n>>> input2 = torch.rand(1, 4, 5, 5)\n>>> ssim = kornia.losses.SSIM(5, reduction='none')\n>>> loss = ssim(input1, input2) # 1x4x5x5\n\nclass InverseDepthSmoothnessLoss[source]\n\nCriterion that computes image-aware inverse depth smoothness loss.\n\n$\\text{loss} = \\left | \\partial_x d_{ij} \\right | e^{-\\left \\| \\partial_x I_{ij} \\right \\|} + \\left | \\partial_y d_{ij} \\right | e^{-\\left \\| \\partial_y I_{ij} \\right \\|}$\nShape:\n• Inverse Depth: $$(N, 1, H, W)$$\n\n• Image: $$(N, 3, H, W)$$\n\n• Output: scalar\n\nExamples:\n\n>>> idepth = torch.rand(1, 1, 4, 5)\n>>> image = torch.rand(1, 3, 4, 5)\n>>> smooth = kornia.losses.DepthSmoothnessLoss()\n>>> loss = smooth(idepth, image)\n\nclass TotalVariation[source]\n\nComputes the Total Variation according to .\n\nShape:\n• Input: $$(N, C, H, W)$$ or $$(C, H, W)$$ where C = number of classes.\n\n• Output: $$(N,)$$ or $$()$$\n\nExamples\n\n>>> kornia.losses.total_variation(torch.ones(3,4,4)) # tensor(0.)\n>>> tv = kornia.losses.TotalVariation()\n>>> output = tv(torch.ones(2,3,4,4)) # tensor([0., 0.])\n>>> output.backward()\n\nReference:\nclass PSNRLoss(max_val: float)[source]\n\nCreates a criterion that calculates the PSNR between 2 images. Given an m x n image, the PSNR is:\n\n$\\text{PSNR} = 10 \\log_{10} \\bigg(\\frac{\\text{MAX}_I^2}{MSE(I,T)}\\bigg)$\n\nwhere\n\n$\\text{MSE}(I,T) = \\frac{1}{mn}\\sum_{i=0}^{m-1}\\sum_{j=0}^{n-1} [I(i,j) - T(i,j)]^2$\n\nand $$\\text{MAX}_I$$ is the maximum possible input value (e.g for floating point images $$\\text{MAX}_I=1$$).\n\nParameters\n\nmax_val (float) – Maximum value of input\n\nShape:\n• input: $$(*)$$\n\n• approximation: $$(*)$$ same shape as input\n\n• output: $$()$$ a scalar\n\nExamples\n\n>>> kornia.losses.psnr_loss(torch.ones(1), 1.2*torch.ones(1), 2)\ntensor(20.0000) # 10 * log(4/((1.2-1)**2)) / log(10)\n\nReference:\n\nhttps://en.wikipedia.org/wiki/Peak_signal-to-noise_ratio#Definition" ]
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https://physics.stackexchange.com/questions/tagged/magnetic-monopoles?tab=newest&page=4
[ "# Questions tagged [magnetic-monopoles]\n\nA magnetic monopole is a hypothetical particle with only one magnetic pole. Their magnetic fields would not be divergence-less. Predicted by certain modern theories, including string theory, supergravity, and various popular grand unified theories.\n\n270 questions\nFilter by\nSorted by\nTagged with\n90 views\n\n### Quantising the magnetic monopoles the make Maxwell symmetric\n\nI don't believe this has already been asked, but I might be wrong; sorry. One can add a magnetic charge density/magnetic monopoles to Maxwell's equations to make the theory symmetric between the ...\n81 views\n\n### Would magnetic charge (i.e., magnetic monopoles) be Lorentz invariant if it existed?\n\nWould magnetic charge (i.e., magnetic monopoles) be Lorentz invariant if it existed? It is clear that Maxwell's equations in themselves permit magnetic charges but what would their relativistic ...\n619 views\n\n### Total angular momentum of electron in a magnetic field\n\nIn this question: Electron in the proximity of a magnetic monopole It is stated that for an electron in the magnetic field of a monopole, $\\vec{B}(\\vec{r}) = \\frac{g}{r^3}\\vec{r}$ , that the ...\n82 views\n\n### What is the assurance that a magnetic monopole (if created) can be detected through the MoEDAL Experiment?\n\nPlease refer to this paper (http://web.mit.edu/ellisk/www/monopoles.pdf). It says in the Section B, ... One search technique relies upon the fact that, just as an electric current produces a ...\n67 views\n\n### What will be the motion of an electron kept at rest in the vicinity of a magnetic monopole?\n\nConsider that both an electron and a magnetic monopole are kept at rest close to each other. Will this result any motion of the electron? If yes, please give a derivation of the equation of motion.\n840 views\n\n### Magnetic monopoles and special relativity\n\nI was thinking about magnetism as a product of special relativity and the result of this approach to the magnetic monopoles. So if magnetism is a product of electricity(like electricity from another ...\n109 views\n\n### Quantum mechanical monopole [duplicate]\n\nRecently in physics news, scientists have experimentally discovered the so-called quantum mechanical monopole. It seems that a quantum mechanical monopole is different from a magnetic monopole. So my ...\n265 views\n\n### Poincaré' lemma and EM potential $A^{\\mu}$\n\nMy lecturer said that given the sourceless Maxwell's equations $$\\partial_{\\mu}\\, ^ *F^{\\mu\\nu} = 0$$, we can find a solution $$F^{\\mu\\nu} = \\partial_{\\mu}A_{\\nu} - \\partial_{\\nu}A_{\\mu},$$ that ...\n87 views\n\n### Is this a possible monopole setup, or will it cancel out? [duplicate]\n\nSuppose you have a ball that is covered in magnets, in which the North Pole of every magnet faces outward. Is this a monopole? Or at least \"functions\" as a monopole? And what will happen if you placed ...\n3k views\n\n### What do we mean with magnetic monopole and dipole?\n\nWhat do we mean with magnetic monopole and dipole? I can not find a way to relate magnetic monopoles and dipoles with electric ones. I do not understand their outcomes. Also,what is their role in ...\n131 views\n\n### conducting hollow sphere in magnetic monopole\n\nif a hollow copper sphere(or any conducting hollow sphere) is connected to dc at points diametrical and a magnetic monopole is right at the center of the sphere then will there be any movement of the ...\n640 views\n\n### Is the Lorentz force expression valid for magnetic field created by a magnetic monopole?\n\nWill the Lorentz force expression be valid for a magnetic field created by a magnetic monopole? I haven't seen any derivation of Lorentz force expression yet and I don't know whether it was derived ...\n166 views\n\n### Experiments or phenomenon to falsify the existence of monopoles?\n\nThis is a related, but opposite question to this one I have heard about a lot of things regarding elementary and GUT magnetic monopoles, as well the quasiparticle monopoles in spin ice Since there's ...\n2k views\n\n### Is it possible to make an electromagnet with two like ends?\n\nIs it possible to create an electromagnet with one continuous wire with 2 like poles (i.e. both ends either north or south)?. Visualising it with the right hand screw rule for current carrying coils,...\n1k views\n\n### Does magnetic monopole violate $U(1)$ gauge symmetry?\n\nDoes a magnetic monopole violate $U(1)$ gauge symmetry? In what sense and why? Insofar as I know, there are at least two types of magnetic monopoles. One is the Dirac monopole while the other is the ...\n11k views\n\n### How to demagnetize a magnet temporarily?\n\nHow to demagnetize a magnet temporarily? Is it possible making a permanent magnet demagnetize for a short span of time? Or at least long time but temporarily!\n672 views\n\n### How will SR EM Lagrangian change if we find a magnetic charge?\n\nWhen we introduce electromagnetic field in Special Relativity, we add a term of $$-\\frac e c A_idx^i$$ into Lagrangian. When we then derive equations of motion, we get the magnetic field that is ...\n3k views\n\n### How are the Lorentz force, Maxwell's third law and Faraday's law of induction clasically related?\n\nFaraday's law of induction can be used in any situation where the magnetic flux is changing through a closed conducting loop. While giving the correct answer, it seems to me that for the following ...\n450 views\n\n### Need of vector potential in quantum mechanics\n\nI need your opinions. Why is the vector potential of a magnetic field important (or even necessary) to quantum mechanics? Why it has to be defined everywhere? Is there any fundamental reason you can ...\n52 views\n\n### Vector potential in presence of monopole [duplicate]\n\nIn this paper http://www.hcs.harvard.edu/~jus/0302/song.pdf when Song was explaining dirac string. He said \"In the presence of a magnetic monopole, the vector potential cannot be defined everywhere. ...\n1k views\n\n### Why do Magnetic North and South Poles never exist by themselves? [duplicate]\n\nAs the title suggest, whilst I was reading I saw written that a 'Magnetic North or South Pole has never been found by itself'. And I was wondering why this was?\n1k views\n\n543 views\n\n### No monopoles in the Weinberg-Salam model\n\nI'm reading Chapter 10.4 on the 't Hooft-Polyakov monopoles in Ryder's Quantum Field Theory. On page 412 he explains why magnetic monopoles cannot appear in the Weinberg-Salam model. I'm I right by ...\n309 views\n\n### A graphical proof that the $SU(2)/\\mathbb{Z}_2$ vortex is non-orientable\n\nThe text, see , compares the vortex solutions of a spontaneously broken symmetry $U(1) \\rightarrow 1$ and $SU(2)\\rightarrow U(1) \\rightarrow \\mathbb{Z}_2$. The vortices can be classified by ...\n23 views\n\n### quasiparticles resembling monopole energy generator\n\nIs it possible using quasiparticles that superficially resemble monopoles to make a generator providing a sustainable source of power\n183 views\n\n### Magnetic monopole\n\nI've been watching a lot of videos claiming to make monopoles either by dropping the temperature to 5k (-268° C) and splitting the magnet or by pouring molten metal through an induction system and was ..." ]
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https://www.geeksforgeeks.org/move-last-element-to-front-of-a-given-linked-list/?ref=rp
[ "Related Articles\n\n# Move last element to front of a given Linked List\n\n• Difficulty Level : Basic\n• Last Updated : 30 Oct, 2020\n\nWrite a function that moves the last element to the front in a given Singly Linked List. For example, if the given Linked List is 1->2->3->4->5, then the function should change the list to 5->1->2->3->4.\n\nAlgorithm:\nTraverse the list till last node. Use two pointers: one to store the address of last node and other for address of second last node. After the end of loop do following operations.\ni) Make second last as last (secLast->next = NULL).\n\n## C++\n\n `/* CPP Program to move last element ``to front in a given linked list */``#include ``using` `namespace` `std;`` ` `/* A linked list node */``class` `Node ``{ ``    ``public``:``    ``int` `data; ``    ``Node *next; ``}; `` ` `/* We are using a double pointer``head_ref here because we change ``head of the linked list inside ``this function.*/``void` `moveToFront(Node **head_ref) ``{ ``    ``/* If linked list is empty, or ``    ``it contains only one node, ``    ``then nothing needs to be done,``    ``simply return */``    ``if` `(*head_ref == NULL || (*head_ref)->next == NULL) ``        ``return``; `` ` `    ``/* Initialize second last``    ``and last pointers */``    ``Node *secLast = NULL; ``    ``Node *last = *head_ref; `` ` `    ``/*After this loop secLast contains``    ``address of second last node and ``    ``last contains address of last node in Linked List */``    ``while` `(last->next != NULL) ``    ``{ ``        ``secLast = last; ``        ``last = last->next; ``    ``} `` ` `    ``/* Set the next of second last as NULL */``    ``secLast->next = NULL; `` ` `    ``/* Set next of last as head node */``    ``last->next = *head_ref; `` ` `    ``/* Change the head pointer``    ``to point to last node now */``    ``*head_ref = last; ``} `` ` `/* UTILITY FUNCTIONS */``/* Function to add a node ``at the beginning of Linked List */``void` `push(Node** head_ref, ``int` `new_data) ``{ ``    ``/* allocate node */``    ``Node* new_node = ``new` `Node();`` ` `    ``/* put in the data */``    ``new_node->data = new_data; `` ` `    ``/* link the old list off the new node */``    ``new_node->next = (*head_ref); `` ` `    ``/* move the head to point to the new node */``    ``(*head_ref) = new_node; ``} `` ` ` ` `/* Function to print nodes in a given linked list */``void` `printList(Node *node) ``{ ``    ``while``(node != NULL) ``    ``{ ``        ``cout << node->data << ``\" \"``; ``        ``node = node->next; ``    ``} ``} `` ` `/* Driver code */``int` `main() ``{ ``    ``Node *start = NULL; `` ` `    ``/* The constructed linked list is: ``    ``1->2->3->4->5 */``    ``push(&start, 5); ``    ``push(&start, 4); ``    ``push(&start, 3); ``    ``push(&start, 2); ``    ``push(&start, 1); `` ` `    ``cout<<``\"Linked list before moving last to front\\n\"``; ``    ``printList(start); `` ` `    ``moveToFront(&start); `` ` `    ``cout<<``\"\\nLinked list after removing last to front\\n\"``; ``    ``printList(start); `` ` `    ``return` `0; ``} `` ` `// This code is contributed by rathbhupendra`\n\n## C\n\n `/* C Program to move last element to front in a given linked list */``#include``#include`` ` `/* A linked list node */``struct` `Node``{``    ``int` `data;``    ``struct` `Node *next;``};`` ` `/* We are using a double pointer head_ref here because we change``   ``head of the linked list inside this function.*/``void` `moveToFront(``struct` `Node **head_ref)``{``    ``/* If linked list is empty, or it contains only one node,``      ``then nothing needs to be done, simply return */``    ``if` `(*head_ref == NULL || (*head_ref)->next == NULL)``        ``return``;`` ` `    ``/* Initialize second last and last pointers */``    ``struct` `Node *secLast = NULL;``    ``struct` `Node *last = *head_ref;`` ` `    ``/*After this loop secLast contains address of second last``    ``node and last contains address of last node in Linked List */``    ``while` `(last->next != NULL)``    ``{``        ``secLast = last;``        ``last = last->next;``    ``}`` ` `    ``/* Set the next of second last as NULL */``    ``secLast->next = NULL;`` ` `    ``/* Set next of last as head node */``    ``last->next = *head_ref;`` ` `    ``/* Change the head pointer to point to last node now */``    ``*head_ref = last;``}`` ` `/* UTILITY FUNCTIONS */``/* Function to add a node at the beginning of Linked List */``void` `push(``struct` `Node** head_ref, ``int` `new_data)``{``    ``/* allocate node */``    ``struct` `Node* new_node =``        ``(``struct` `Node*) ``malloc``(``sizeof``(``struct` `Node));`` ` `    ``/* put in the data  */``    ``new_node->data  = new_data;`` ` `    ``/* link the old list off the new node */``    ``new_node->next = (*head_ref);`` ` `    ``/* move the head to point to the new node */``    ``(*head_ref)    = new_node;``}`` ` ` ` `/* Function to print nodes in a given linked list */``void` `printList(``struct` `Node *node)``{``    ``while``(node != NULL)``    ``{``        ``printf``(``\"%d \"``, node->data);``        ``node = node->next;``    ``}``}`` ` `/* Driver program to test above function */``int` `main()``{``    ``struct` `Node *start = NULL;`` ` `    ``/* The constructed linked list is:``     ``1->2->3->4->5 */``    ``push(&start, 5);``    ``push(&start, 4);``    ``push(&start, 3);``    ``push(&start, 2);``    ``push(&start, 1);`` ` `    ``printf``(``\"\\n Linked list before moving last to front\\n\"``);``    ``printList(start);`` ` `    ``moveToFront(&start);`` ` `    ``printf``(``\"\\n Linked list after removing last to front\\n\"``);``    ``printList(start);`` ` `    ``return` `0;``}`\n\n## Java\n\n `/* Java Program to move last element to front in a given linked list */``class` `LinkedList``{``    ``Node head;  ``// head of list``  ` `    ``/* Linked list Node*/``    ``class` `Node``    ``{``        ``int` `data;``        ``Node next;``        ``Node(``int` `d) {data = d; next = ``null``; }``    ``}`` ` `    ``void` `moveToFront()``    ``{``        ``/* If linked list is empty or it contains only``           ``one node then simply return. */``           ``if``(head == ``null` `|| head.next == ``null``) ``              ``return``;`` ` `        ``/* Initialize second last and last pointers */``        ``Node secLast = ``null``;``        ``Node last = head;`` ` `        ``/* After this loop secLast contains address of ``           ``second last  node and last contains address of ``           ``last node in Linked List */``        ``while` `(last.next != ``null``)  ``        ``{``           ``secLast = last;``           ``last = last.next; ``        ``}`` ` `        ``/* Set the next of second last as null */``        ``secLast.next = ``null``;`` ` `        ``/* Set the next of last as head */``        ``last.next = head;`` ` `        ``/* Change head to point to last node. */``        ``head = last;``    ``}                 `` ` `                     ` `    ``/* Utility functions */`` ` `    ``/* Inserts a new Node at front of the list. */``    ``public` `void` `push(``int` `new_data)``    ``{``        ``/* 1 & 2: Allocate the Node &``                  ``Put in the data*/``        ``Node new_node = ``new` `Node(new_data);``  ` `        ``/* 3. Make next of new Node as head */``        ``new_node.next = head;``  ` `        ``/* 4. Move the head to point to new Node */``        ``head = new_node;``    ``}`` ` `    ``/* Function to print linked list */``    ``void` `printList()``    ``{``        ``Node temp = head;``        ``while``(temp != ``null``)``        ``{``           ``System.out.print(temp.data+``\" \"``);``           ``temp = temp.next;``        ``}  ``        ``System.out.println();``    ``}`` ` `     ``/* Driver program to test above functions */``    ``public` `static` `void` `main(String args[])``    ``{``        ``LinkedList llist = ``new` `LinkedList();``        ``/* Constructed Linked List is 1->2->3->4->5->null */``        ``llist.push(``5``);``        ``llist.push(``4``);``        ``llist.push(``3``);``        ``llist.push(``2``);``        ``llist.push(``1``);``         ` `        ``System.out.println(``\"Linked List before moving last to front \"``);``        ``llist.printList();``         ` `        ``llist.moveToFront();``         ` `        ``System.out.println(``\"Linked List after moving last to front \"``);``        ``llist.printList();``    ``}``} ``/* This code is contributed by Rajat Mishra */`\n\n## Python3\n\n `# Python3 code to move the last item to front``class` `Node:``    ``def` `__init__(``self``, data):``        ``self``.data ``=` `data``        ``self``.``next` `=` `None`` ` `class` `LinkedList:``    ``def` `__init__(``self``):``        ``self``.head ``=` `None`` ` `    ``# Function to add a node ``    ``# at the beginning of Linked List``    ``def` `push(``self``, data):``        ``new_node ``=` `Node(data)``        ``new_node.``next` `=` `self``.head``        ``self``.head ``=` `new_node``         ` `    ``# Function to print nodes in a``    ``# given linked list``    ``def` `printList(``self``):``        ``tmp ``=` `self``.head``        ``while` `tmp ``is` `not` `None``:``            ``print``(tmp.data, end``=``\", \"``)``            ``tmp ``=` `tmp.``next``        ``print``()`` ` `    ``# Function to bring the last node to the front``    ``def` `moveToFront(``self``):``        ``tmp ``=` `self``.head``        ``sec_last ``=` `None` `# To maintain the track of``                        ``# the second last node`` ` `    ``# To check whether we have not received ``    ``# the empty list or list with a single node``        ``if` `not` `tmp ``or` `not` `tmp.``next``: ``            ``return`` ` `        ``# Iterate till the end to get``        ``# the last and second last node ``        ``while` `tmp ``and` `tmp.``next` `:``            ``sec_last ``=` `tmp``            ``tmp ``=` `tmp.``next`` ` `        ``# point the next of the second``        ``# last node to None``        ``sec_last.``next` `=` `None`` ` `        ``# Make the last node as the first Node``        ``tmp.``next` `=` `self``.head``        ``self``.head ``=` `tmp`` ` `# Driver Code``if` `__name__ ``=``=` `'__main__'``:``    ``llist ``=` `LinkedList()``     ` `    ``# swap the 2 nodes``    ``llist.push(``5``)``    ``llist.push(``4``)``    ``llist.push(``3``)``    ``llist.push(``2``)``    ``llist.push(``1``)``    ``print` `(``\"Linked List before moving last to front \"``)``    ``llist.printList()``    ``llist.moveToFront()``    ``print` `(``\"Linked List after moving last to front \"``)``    ``llist.printList()`\n\n## C#\n\n `/* C# Program to move last element to front in a given linked list */``using` `System;``class` `LinkedList ``{ ``    ``Node head; ``// head of list `` ` `    ``/* Linked list Node*/``    ``public` `class` `Node ``    ``{ ``        ``public` `int` `data; ``        ``public` `Node next; ``        ``public` `Node(``int` `d) {data = d; next = ``null``; } ``    ``} `` ` `    ``void` `moveToFront() ``    ``{ ``        ``/* If linked list is empty or it contains only ``        ``one node then simply return. */``        ``if``(head == ``null` `|| head.next == ``null``) ``            ``return``; `` ` `        ``/* Initialize second last and last pointers */``        ``Node secLast = ``null``; ``        ``Node last = head; `` ` `        ``/* After this loop secLast contains address of ``        ``second last node and last contains address of ``        ``last node in Linked List */``        ``while` `(last.next != ``null``) ``        ``{ ``        ``secLast = last; ``        ``last = last.next; ``        ``} `` ` `        ``/* Set the next of second last as null */``        ``secLast.next = ``null``; `` ` `        ``/* Set the next of last as head */``        ``last.next = head; `` ` `        ``/* Change head to point to last node. */``        ``head = last; ``    ``}                 `` ` `                     ` `    ``/* Utility functions */`` ` `    ``/* Inserts a new Node at front of the list. */``    ``public` `void` `push(``int` `new_data) ``    ``{ ``        ``/* 1 & 2: Allocate the Node & ``                ``Put in the data*/``        ``Node new_node = ``new` `Node(new_data); `` ` `        ``/* 3. Make next of new Node as head */``        ``new_node.next = head; `` ` `        ``/* 4. Move the head to point to new Node */``        ``head = new_node; ``    ``} `` ` `    ``/* Function to print linked list */``    ``void` `printList() ``    ``{ ``        ``Node temp = head; ``        ``while``(temp != ``null``) ``        ``{ ``        ``Console.Write(temp.data+``\" \"``); ``        ``temp = temp.next; ``        ``} ``        ``Console.WriteLine(); ``    ``} `` ` `    ``/* Driver program to test above functions */``    ``public` `static` `void` `Main(String []args) ``    ``{ ``        ``LinkedList llist = ``new` `LinkedList(); ``        ``/* Constructed Linked List is 1->2->3->4->5->null */``        ``llist.push(5); ``        ``llist.push(4); ``        ``llist.push(3); ``        ``llist.push(2); ``        ``llist.push(1); ``         ` `        ``Console.WriteLine(``\"Linked List before moving last to front \"``); ``        ``llist.printList(); ``         ` `        ``llist.moveToFront(); ``         ` `        ``Console.WriteLine(``\"Linked List after moving last to front \"``); ``        ``llist.printList(); ``    ``} ``} ``// This code is contributed by Arnab Kundu`\n\nOutput:\n``` Linked list before moving last to front\n1 2 3 4 5\nLinked list after removing last to front\n5 1 2 3 4```\n\nTime Complexity: O(n) where n is the number of nodes in the given Linked List.\n\nPlease write comments if you find any bug in the above code/algorithm, or find other ways to solve the same problem.\n\nAttention reader! Don’t stop learning now. Get hold of all the important DSA concepts with the DSA Self Paced Course at a student-friendly price and become industry ready.  To complete your preparation from learning a language to DS Algo and many more,  please refer Complete Interview Preparation Course.\n\nIn case you wish to attend live classes with experts, please refer DSA Live Classes for Working Professionals and Competitive Programming Live for Students.\n\nMy Personal Notes arrow_drop_up" ]
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https://www.hindawi.com/journals/aaa/2010/902638/
[ "/ / Article\n\nResearch Article | Open Access\n\nVolume 2010 |Article ID 902638 | https://doi.org/10.1155/2010/902638\n\nValery Serov, \"Green's Function and Convergence of Fourier Series for Elliptic Differential Operators with Potential from Kato Space\", Abstract and Applied Analysis, vol. 2010, Article ID 902638, 18 pages, 2010. https://doi.org/10.1155/2010/902638\n\n# Green's Function and Convergence of Fourier Series for Elliptic Differential Operators with Potential from Kato Space\n\nAccepted04 Feb 2010\nPublished02 Mar 2010\n\n#### Abstract\n\nWe consider the Friedrichs self-adjoint extension for a differential operator of the form , which is defined on a bounded domain , (for we assume that is a finite interval). Here is a formally self-adjoint and a uniformly elliptic differential operator of order 2m with bounded smooth coefficients and a potential is a real-valued integrable function satisfying the generalized Kato condition. Under these assumptions for the coefficients of and for positive large enough we obtain the existence of Green's function for the operator and its estimates up to the boundary of . These estimates allow us to prove the absolute and uniform convergence up to the boundary of of Fourier series in eigenfunctions of this operator. In particular, these results can be applied for the basis of the Fourier method which is usually used in practice for solving some equations of mathematical physics.\n\n#### 1. Introduction\n\nLet be a bounded domain in () with smooth boundary. We consider on an elliptic differential operator of the form\n\nwhere\n\nis a formally self-adjoint differential operator of order . Here with . The coefficients of the operator are assumed to be the complex-valued (in general) bounded smooth functions on the domain for all such that are real valued for and this operator satisfies the uniform ellipticity condition\n\nwith some constant , for all and all . We assume that the potential is a real-valued -function satisfying the generalized Kato condition, that is,\n\nwhere function for is defined by\n\nFor we denote by the -based Sobolev space, where indicates the “degree” of the smoothness; by we denote the closure of in . We denote also by Besov space (where indicates the smoothness) with the same notation for as for Sobolev space. The definition of Sobolev and Besov spaces as well as the embedding theorems for these spaces can be found in [1, 2].\n\nDue to (1.4)-(1.5) the function satisfies all conditions of Theorem of with and therefore for any we have the following inequality:\n\nwhere the constant depends only on and , the constant depends only on , and , and the value is defined by\n\nwhere is as in (1.5).\n\nSince the domain is bounded then tends to as . It immediately implies that there is a constant such that\n\nfor all . Since is positive for sufficiently large it has a positive self-adjoint Friedrichs extension such that\n\nWe define the Friedrichs extension of to be such that\n\nThe domain of is given by\n\nIt is also well known that this extension has a purely discrete spectrum of finite multiplicity having the only one accumulation point at infinity () and a complete orthonormal system of eigenfunctions in .\n\nTo each function we can assign the formal series\n\nwhere are the Fourier coefficients of with respect to the system .\n\nThe study of elliptic differential operators with smooth coefficients on a bounded domain with smooth boundary has a long history. We restrict the bibliographical remarks to the works that are of interest from the viewpoint of the present article.\n\nThe estimates for the Green's function and convergence of spectral expansions of a general elliptic differential operator of order with smooth coefficients on a bounded domain have been studied by many authors. We refer to a four-volume monograph of Hörmander [4, 5], the works of Alimov , Gårding , Krasovskiĭ [11, 12], Schechter and others. We mention also the papers of the author of the present which deal with the operators whose coefficients may have local singularities of specific order on an arbitrary smooth surface whose dimension is strictly less than that of the original domain. As to elliptic operators of order whose coefficients may have singularities in , similar results have been mainly obtained for the Schrödinger operators on with from or in any dimensions but with which may have given singularity at one point. For such results, see Alimov and Joó , Ashurov , Ashurov and Faiziev , Khalmukhamedov [19, 20], Serov [21, 22], Serov and Buzurnyuk , and others.Some survey of resent results concerning theory of elliptic differential operators of order can be found in the articles of Davies [24, 25].\n\nThe aim of this paper is to prove the following results.\n\nTheorem 1.1. Suppose that satisfies condition (1.4), then there exist constants , and such that for all the Green function of the operator exists and satisfies the following estimates:for all and .\n\nWithout loss of generality, in the following theorem we assume that is positive.\n\nTheorem 1.2. The Fourier series (1.12) converges absolutely and uniformly on the domain for any function from the domain of the operator for .\n\nOne of the main results of the present paper is Theorem 1.1 which concerns the estimates up to the boundary of the domain for the Green's function of an elliptic differential operator of order with singular potential from the generalized Kato space. In all previous publications, as far as we know, the estimates for the Green function are proved on an arbitrary compact subset from the domain and for the case when the coefficients of operator are either smooth or have some special type of singularities.\n\nAnother main result of this paper is Theorem 1.2. It gives some sufficient conditions which provide the absolute convergence up to the boundary of of Fourier series in eigenfunctions for functions from the domain of some power of this operator. In addition to Theorem 1.2, we would like to take into consideration Theorem 3.7 (see Section 3 of the paper) which is the generalization of the well-known result of Peetre (see ) to the operators with singular coefficients. It can be mentioned also here that in the scale of the spaces associated with some powers of our operator the results of Theorems 1.2 and 3.7 are sharp (see, e.g., ).\n\nThis paper is organized such that Theorem 1.1 is proved in Section 2 and Theorem 1.2 in Section 3. Some additional theorems about the absolute convergence of Fourier series are also proved in Section 3.\n\n#### 2. Green's Function\n\nIn this section we obtain the estimates for the Green's function of the operator when is positive and sufficiently large.\n\nDefinition 2.1. For and , a locally integrable function on is called a fundamental solution for an operator if and only if\n\nEquation (2.1) holds in the sense of distributions, that is,\n\nfor all , where\n\nis the transpose of .\n\nWe will use the following result.\n\nProposition 2.2. There exists such that for any , the differential operator has a fundamental solution . Furthermore, for any multi-index , , there are constants , such that the following estimates hold: for all and .\n\nThe proof of Proposition 2.2 can be found in .\n\nWe will look for the fundamental solution of the operator , for positive and large enough, as a solution of the integral equation\n\nwhere is the fundamental solution of the operator . By Proposition 2.2, exists and belongs (at least) to uniformly with respect to from .\n\nWe need the following lemma, which may have interest of its own right.\n\nLemma 2.3. Assume that satisfies condition (1.4), then there is such that for all the fundamental solution exists as a solution of the integral equation (2.4) and satisfies the following estimates: with some positive constant , where is as in Proposition 2.2 and .\n\nProof. We solve the integral equation (2.4) by iterations. For any , we denote We will prove by induction that there is such that for all and for each where , is as in the Proposition 2.2 and is defined as It is clear that for estimate (2.6) holds. And it is also clear that (2.6) holds for the case when for each by choosing large enough.\nIn the case (considering two possibilities and ) in order to prove (2.6) it is enough to prove that there exists such that for all and , where constant is as in Proposition 2.2.\nIndeed, since for we have where is as in (1.4), then we can estimate the left-hand side of (2.8) as follows. If then the integrals in the latter equality tend to zero as . The reason is due to that condition (1.4) the measure of the set tends to zero as . If then this integral can be estimated by where some positive constant depends only on and dimension .\nIn the case (considering two possibility and ) it can be proved that Then, instead of estimate (2.10) we obtain\nCombining these two facts (including (2.10) and (2.12)), (2.8) and choosing appropriately (with respect to ) we may conclude that inequality (2.6) is proved. Since the solution of the integral equation (2.4) is given by the series the estimates (2.6) prove also Lemma 2.3.\n\nAs a consequence of Lemma 2.3 and Proposition 2.2 we can obtain the estimates for the derivatives of order of the fundamental solution . For the derivatives of we use the following representation:\n\nThe following corollary holds.\n\nCorollary 2.4. Assume that satisfies the condition (1.4), then for the derivatives of the fundamental solution of order the following estimates hold: for some constant , for all and (where and are as in Lemma 2.3).\n\nThe proof of the corollary follows immediately from the integral representation for the derivatives of , estimates for the derivatives of in Proposition 2.2, estimates for in Lemma 2.3, and estimates for the kernels with weak singularities.\n\nLet us note that the fundamental solution , which was obtained in Lemma 2.3, belongs (at least) to uniformly with respect to from .\n\nNow we are in the position to introduce the Green's function of the operator . If is sufficiently large then the operator is positive and its inverse\n\nis a bounded operator. It is also an integral operator with kernel denoted by . If we use for this integral operator the symbol then we have\n\nDefinition 2.5. The kernel of the integral operator is called the Green's function of the operator .\n\nProof of Theorem 1.1. For , let and be compact sets, each of them having a smooth boundary, with such that Here denotes the distance between the sets and .\nLet be a fundamental solution of the operator for and sufficiently large. We choose the function such that and set By this equation the function is well defined for all . Clearly, for , . We will show that is a parametrix for . To prove this, let us introduce the function and corresponding integral operator with kernel where and are integral operators in with kernels and , respectively. Then it follows from (2.20) that where If we denote by the kernel of the integral operator , then it follows from (2.24) that for any , and the kernel (see (2.20)) has the form where . As a matter of fact we cannot characterize and estimate the kernel from (2.22)—(2.24). That is why we will proceed a little bit differently, as follows. Equality (2.19) implies that in the sense of distributions the following representation holds: where ( is considered here as a parameter) and sufficiently large. The function in (2.27) will be of the form with the differential operator having the symbol . It is the polynomial in of order and therefore the differential operators are of order . This fact allows us to estimate the function (in comparison with ). Indeed, by the choice of , only on the set and therefore the representation (2.28) and Corollary 2.4 imply that the following estimate holds: for all and with as in Corollary 2.4.\n\nNow we need the following lemma.\n\nLemma 2.6. For all where is as in (2.28) and is defined as in (2.19).\n\nProof. We can rewrite (2.27) in the operator form as or (using (2.20)) The latter equation implies and therefore (using (2.20) again) But this is equivalent to (2.30). Thus, this lemma is proved.\n\nIn order to finish the proof of Theorem 1.1 let us introduce new functions and which are obtained from and multiplying by\n\nrespectively, where is as in Corollary 2.4. Then (2.30) (see Lemma 2.6) and estimate (2.29) formally yield the following estimate (for simplicity let us consider here only the case , the cases can be considered similarly)\n\nConsidering two possibilities and the value in the latter brackets can be estimated from above by\n\nThis estimate allows us to get from (2.36) that\n\nSince\n\nthen for large enough (2.38) yields\n\nThus, Theorem 1.1 is completely proved.\n\n#### 3. Convergence of Fourier Series\n\nWithout loss of generality, we assume in this section that is positive. Then by the J. von Neumann spectral theorem for , where with as in Theorem 1.1, the following representation holds:\n\nwhere is real and is the spectral resolution corresponding to the self-adjoint operator . The domain of the operator (3.1) is defined by\n\nIn our case (in the case of pure discrete spectrum), the spectral projector has the form\n\nwhere are the Fourier coefficients of with respect to the system . Hence relations (3.1) and (3.2) become\n\nIn addition, we need a special representation of the negative fractional powers of . If we assume that then using well-known properties of Euler beta-function, one can obtain\n\nwhere is the integral operator with kernel from Section 2. This representation shows that operator (3.6) is also integralwith kernel denoted by . Using Theorem 1.1 of present work and well-known technique (see, e.g., ) it is not so difficult to prove the following estimates\n\nwhere , is as in Theorem 1.1 and the constant depends on .\n\nThe following main lemma holds.\n\nLemma 3.1. For any function and for where with as in Theorem 1.1.\n\nProof. For any , we write and represent the operator as the product Then, by applying the estimates (3.7) and Lemma in , we arrive at (3.8). This completes the proof.\n\nCorollary 3.2. Assume that . There is a constant depending only on , such that the estimate holds uniformly in and .\n\nProof. By the spectral theorem and relation (3.4), we can rewrite inequality (3.8) in the form where are the Fourier coefficients of with respect to the system . Now inequality (3.11) follows by duality. The proof is complete.\n\nRemark 3.3. The inequality (3.11) has an independent interest since it gives the \"bundle\" estimate of the eigenfunctions in the form which holds uniformly in and large enough. Indeed, from (3.12) we have uniformly in . If we chose now then one can immediately obtain (3.13).\n\nNow we are ready to prove Theorem 1.2.\n\nProof of Theorem 1.2. Using the representation (3.4), the inequality (3.11), and the Cauchy-Schwarz-Bunyakovsk 2 inequality, we obtain uniformly in and for any fixed . Now the desired assertion follows from (3.5). Theorem 1.2 is completely proved.\n\nThe estimate (3.13) allows us to obtain a bit more precise result than in Theorem 1.2. Namely, the following corollary holds.\n\nCorollary 3.4. Assume that the function satisfies the condition where are the Fourier coefficients of with respect to the system , then the Fourier series (1.12) converges absolutely and uniformly on .\n\nLet us assume now that the potential satisfies the conditions\n\nthen it is not so difficult to see that for the case conditions (3.17) imply the condition (1.4). For the case this condition (3.17) for must be considered in addition to the condition (1.4). The following result is valid.\n\nTheorem 3.5. Suppose that the potential satisfies conditions (3.17), then for any function , where is given by (3.3).\n\nProof. Using the Sobolev embedding theorem we easily conclude that conditions (3.17) imply the following inclusion: And for any the following inequality holds: Moreover, we may assert that the operator is invertible for large enough. Indeed, since the function belongs to , we have the representation for where denotes the Friedrichs self-adjoint extension for in . Using again the Sobolev embedding theorem and conditions (3.17) we may conclude that the functions and belong to . The results of yield that the operator is invertible with small norm for its inverse operator (if is large enough). This fact and the latter equality imply that for large enough the operator is also invertible and for any we have the following inequality:\nNow let . Then (3.23) implies where belongs to and the convergence to zero in the last term follows from the J. von Neumann spectral theorem. The proof is complete.\n\nThe Sobolev embedding theorem gives the immediate corollary.\n\nCorollary 3.6. Let . Then for any with holds uniformly in .\n\nThe next theorem gives us some sufficient conditions which provide the absolute and uniform convergence of Fourier series (1.12) in the classical Besov and Sobolev spaces. Following , we use the symbol .\n\nTheorem 3.7. Assume that the potential belongs to Sobolev space , where is an entire part of , then for any function from Besov space the Fourier series (1.12) converges absolutely and uniformly on the domain .\n\nProof. Using the Sobolev embedding theorem and the following representation: where are some constants, we can conclude that the condition for the potential implies the following inclusion: Then using the results of (see Theorem ) we may conclude that Consequently, by Theorem of and by Peetre's method of real interpolation (see, e.g., ), we have But the latter space is the interpolation space of Peetre (see ) for the elliptic differential operator of order . Since estimate (3.13) for the spectral function holds in our case, we can apply the results of and conclude that the proof of this theorem is complete.\n\nRemark 3.8. If is even then the statement of this theorem holds for any function from Besov space due to the equality (see Theorem of ) And the Sobolev embedding theorem for Besov spaces (see, e.g., ) implies that this theorem also holds for any function from Sobolev space with and arbitrary integer .\n\n#### Acknowledgment\n\nThis work was supported by the Academy of Finland (application no. 213476, Finnish Programme for Centres of Excellence in Research 2006–2011).\n\n1. R. A. Adams and J. F. Fournier, Sobolev Spaces, Academic Press, New York, NY, USA, 2nd edition, 2003.\n2. H. Triebel, Interpolation Theory Function Spaces Differential Operators, Mir, Moscow, Russia, 1980.\n3. M. Schechter, Spectra of Partial Differential Operators, vol. 1, North-Holland, Amsterdam, The Netherlands, 1971, North-Holland Series in Applied Mathematics and Mechanics. View at: MathSciNet\n4. L. Hörmander, The Analysis of Linear Partial Differential Equations, Vols. 1-2, Springer, New York, NY, USA, 1983.\n5. L. Hörmander, The Analysis of Linear Partial Differential Equations, Vols. 3-4, Springer, New York, NY, USA, 1985.\n6. S. A. Alimov, “Fractional powers of elliptic operators and isomorphism of classes of differentiable functions,” Differentsial'nye Uravneniya, vol. 8, pp. 1609–1626, 1972 (Russian). View at: Google Scholar | MathSciNet\n7. S. A. Alimov, “Uniform convergence and summability of the spectral expansions of functions from ${L}_{p}^{\\alpha }$,” Differentsial'nye Uravneniya, vol. 9, no. 4, pp. 669–681, 1973. View at: Google Scholar | MathSciNet\n8. S. A. Alimov, “The spectral expansions of functions belonging to ${H}_{p}^{\\alpha }$,” Matematicheskii Sbornik, vol. 101(143), no. 1, pp. 3–20, 1976. View at: Google Scholar | MathSciNet\n9. S. A. Alimov, “On the absolute convergence of spectral expansions,” Doklady Akademii Nauk, vol. 342, no. 4, pp. 446–448, 1995 (Russian). View at: Google Scholar | MathSciNet\n10. L. Gårding, “Dirichlet's problem for linear elliptic partial differential equations,” Mathematica Scandinavica, vol. 1, pp. 55–72, 1953. View at: Google Scholar\n11. J. P. Krasovskiĭ, “Isolation of the singularity in Green's function,” Izvestiya Akademii Nauk SSSR. Seriya Matematicheskaya, vol. 31, pp. 977–1010, 1967 (Russian). View at: Google Scholar | MathSciNet\n12. J. P. Krasovskiĭ, “Properties of Green's functions, and generalized solutions of elliptic boundary value problems,” Doklady Akademii Nauk SSSR, vol. 184, no. 2, pp. 270–273, 1969. View at: Google Scholar | MathSciNet\n13. V. S. Serov, “On the fundamental solution of a differential operator with a singularity,” Differentsial'nye Uravneniya, vol. 23, no. 3, pp. 531–534, 1987.\n14. V. S. Serov, “The absolute convergence of spectral expansions of operators with a singularity,” Differentsial'nye Uravneniya, vol. 28, no. 1, pp. 127–136, 1992. View at: Google Scholar | MathSciNet\n15. V. S. Serov, “On spectral expansions of functions in ${H}_{p}^{\\alpha }$ for a differential operator with a singularity on the surface,” Doklady Rossiĭskaya Akademiya Nauk, vol. 340, no. 1, pp. 26–28, 1995. View at: Google Scholar | MathSciNet\n16. S. A. Alimov and I. Joó, “On the Riesz summability of eigenfunction expansions,” Acta Scientiarum Mathematicarum, vol. 45, no. 1–4, pp. 5–18, 1983.\n17. R. R. Ashurov, “Asymptotic behavior of a spectral function of the Schrödinger operator with potential $q\\in {L}^{2}\\left({ℝ}^{3}\\right)$,” Differentsial'nye Uravneniya, vol. 23, no. 1, pp. 169–172, 1987. View at: Google Scholar | MathSciNet\n18. R. R. Ashurov and J. E. Faiziev, “On eigenfunction expansions associated with the Schrödinger operator with a singular potential,” Differentsial'nye Uravneniya, vol. 41, no. 2, pp. 241–249, 2005. View at: Google Scholar\n19. A. R. Khalmukhamedov, “Eigenfunction expansions for the Schrödinger operator with singular potentials,” Differentsial'nye Uravneniya, vol. 20, no. 9, pp. 1642–1645, 1984. View at: Google Scholar\n20. A. R. Khalmukhamedov, “Convergence of spectral expansion for a singular operator,” Differentsial'nye Uravneniya, vol. 22, no. 12, pp. 2107–2117, 1986. View at: Google Scholar\n21. V. S. Serov, “On the convergence of Fourier series in eigenfunctions of the Schrödinger operator with Kato potential,” Matematicheskie Zametki, vol. 67, no. 5, pp. 755–763, 2000. View at: Publisher Site | Google Scholar | MathSciNet\n22. V. S. Serov, “Fundamental solution and Fourier series in eigenfunctions of degenerate elliptic operator,” Journal of Mathematical Analysis and Applications, vol. 329, no. 1, pp. 132–144, 2007. View at: Publisher Site | Google Scholar | MathSciNet\n23. N. S. Buzurnyuk and V. S. Serov, “On the convergence of Riesz means of spectral expansions that correspond to the Schrödinger operator with a singular potential,” Differentsial'nye Uravneniya, vol. 32, no. 1, pp. 83–89, 1996. View at: Google Scholar | MathSciNet\n24. E. B. Davies, Spectral Theory and Differential Operators, vol. 42 of Cambridge Studies in Advanced Mathematics, Cambridge University Press, Cambridge, UK, 1995. View at: MathSciNet\n25. E. B. Davies, “${L}^{p}$ spectral theory of higher-order elliptic differential operators,” The Bulletin of the London Mathematical Society, vol. 29, no. 5, pp. 513–546, 1997.\n26. J. Peetre, “Absolute convergence of eigenfunction expansions,” Mathematische Annalen, vol. 169, pp. 307–314, 1967.\n\n#### More related articles\n\nArticle of the Year Award: Outstanding research contributions of 2020, as selected by our Chief Editors. Read the winning articles." ]
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http://www.tcs.tifr.res.in/events/finding-optimal-shortest-path-trees?mini=2020-02
[ "# Finding Optimal Shortest Path Trees\n\nSpeaker:\n\n## Time:\n\nFriday, 14 June 2019, 17:15 to 18:15\n\n## Venue:\n\n• A-201 (STCS Seminar Room)\n\n## Organisers:\n\nAbstract: Given a graph $G$ with one source vertex $s$ and several target vertices, a shortest path tree rooted at $s$ is a subgraph of $G$ that preserves distances from $s$ to each of the target vertices. A shortest path tree is called optimal if it has the minimum number of branching vertices (vertices of degree three or more). In this talk, we will see proofs of the following results.\n\n1. Finding an optimal shortest path tree is $\\mathsf{NP}$-hard for general graphs.\n2. Finding an optimal shortest path tree is in $\\mathsf{P}$ for interval graphs.\n\nThis is based on joint work with Jaikumar. A condensed version of this talk will be given at CSR 2019 in Novosibirsk, Russia." ]
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https://www.learn2luvmath.com/post/find-the-equation-of-a-straight-line-given-the-graph
[ "top of page\nSearch\n\n# Find the Equation of a Straight Line given the graph\n\nThe main idea behind finding the equation of a straight line is centered around finding the gradient and y-intercept.\n\nWe have to remember the concept of a gradient is found by finding the change in y divided by the change in x. Its a rate of change of the y values wrt the x values. Some use the method of rise over run too or rise/run. Keep in mind that checking that your answers are correct, especially knowing that when a graph slope right it is positive, and the left will be negative.\n\nThe word intercept means to cut/meet at a point. So y-intercepts meet at the y-axis/cut the y-axis. This is drawn straight from a straight line equation (look for the term without \"x\" or the constant/number)." ]
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https://www.thestudentroom.co.uk/showthread.php?t=73721
[ "Quick binomials qWatch\n\nAnnouncements\nThis discussion is closed.\nThread starter 14 years ago\n#1\nThe coefficeint of x³ in the expanison of (2 + x)(3-ax)4 (to the power of 4)\nis 30. Find the values of the constant a.\n\n0\n14 years ago\n#2\nCoefficient of x^2 in (3 - ax)^4: 4C2 3^2 (-a)^2 = 54a^2\nCoefficient of x^3 in (3 - ax)^4: 4C3 3 (-a)^3 = -12a^3\n\nCoefficient of x^3 in (2 + x)(3 - ax)^4: 2*(-12a^3) + 54a^2 = -24a^3 + 54a^2\n\n30 = -24a^3 + 54a^2\n24a^3 - 54a^2 + 30 = 0\n4a^3 - 9a^2 + 5 = 0\n(a - 1)(4a^2 - 5a - 5) = 0\na = 1 or a = (1/8)(5 +- sqrt(105))\n0\n14 years ago\n#3\n(Original post by jamal1425)\nThe coefficeint of x³ in the expanison of (2 + x)(3-ax)4 (to the power of 4)\nis 30. Find the values of the constant a.\n\n(2+x)(3-ax)^4 = (2+x)(3^4 - 4.3^3.ax + 6.3^2.(ax)^2 - 4.3.(ax)^3 + (ax)^4) = P(x) + (54a^2 -24a^3).x^3\n\nSo 54a^2 - 24a^3 = 30 =>9a^2-4a^3=6\n=>4a^3 - 9a^2 +6 = 0=>(a-1)(4a^2-5a-5)=0\n=>a= 1 or (5 +/- square root of 105)/8\n0\nX\nnew posts", null, "Back\nto top\nLatest\nMy Feed\n\nOops, nobody has postedin the last few hours.\n\nWhy not re-start the conversation?\n\nsee more\n\nSee more of what you like onThe Student Room\n\nYou can personalise what you see on TSR. Tell us a little about yourself to get started.\n\nUniversity open days\n\n• University of Bristol\nWed, 23 Oct '19\n• University of Exeter\nUndergraduate Open Day - Penryn Campus Undergraduate\nWed, 23 Oct '19\n• University of Nottingham\nMini Open Day Undergraduate\nWed, 23 Oct '19\n\nPoll\n\nJoin the discussion\n\nHave you made up your mind on your five uni choices?\n\nYes I know where I'm applying (154)\n59.46%\nNo I haven't decided yet (60)\n23.17%\nYes but I might change my mind (45)\n17.37%" ]
[ null, "https://www.thestudentroom.co.uk/images/v2/icons/arrow_up.svg", null ]
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https://www.b4x.com/android/forum/threads/filter-similar-to-the-vb-net-filter-function.48255/
[ "# Android Code SnippetFilter - (similar to the VB.Net Filter function)\n\nSub Name: Filter\n\nDescription: Returns a String Array containing a subset of the String Array passed as parameter according to the specified filter (Match parameter).\n\nParameters:\nStrings() - the \"source\" String Array\nMatch - the String filter\nInclude - whether to return an array containing Match or that does not contain it\nCaseSensitive - if take into account the difference between upper and lower case\n\nB4X:\n``````Sub Filter(Strings() As String, Match As String, Include As Boolean, CaseSensitive As Boolean) As String()\nDim lstResult As List : lstResult.Initialize\nDim Valid As Boolean\n\nIf Not(Include) Then\nFor i = 0 To Strings.Length - 1\nNext\nEnd If\n\nFor i = 0 To Strings.Length - 1\nValid = False\nIf CaseSensitive Then\nIf Strings(i) = Match Then\nValid = True\nEnd If\nElse\nIf Strings(i).ToLowerCase = Match.ToLowerCase Then\nValid = True\nEnd If\nEnd If\n\nIf Valid Then\nIf Include Then\nElse\nlstResult.RemoveAt(lstResult.IndexOf(Strings(i)))\nEnd If\nEnd If\nNext\n\nDim Result(lstResult.Size) As String\n\nFor i = 0 To lstResult.Size - 1\nResult(i) = lstResult.Get(i)\nNext\n\nReturn Result\n\nEnd Sub``````\n\nExample:\nB4X:\n`````` Dim Strings(5) As String\n\nStrings(0) = \"This\"\nStrings(1) = \"is\"\nStrings(2) = \"just\"\nStrings(3) = \"an\"\nStrings(4) = \"example\"\n\nDim Result() As String = Filter(Strings, \"Is\", True, False)\n' Returns: \"is\"\n\nDim Result() As String = Filter(Strings, \"Is\", True, True)\n' Returns: [empty array]\nLog(\"length= \" & Result.Length)\n\nDim Result() As String = Filter(Strings, \"Is\", False, False)\n' Returns: \"This\" \"just\" \"an\" \"example\"``````\n\nTags: Filter, String, Strings, Texts, Array, Regex\n\nLast edited:\n\nReplies\n24\nViews\n1K\nReplies\n4\nViews\n419\nReplies\n3\nViews\n222\nReplies\n5\nViews\n1K\nReplies\n2\nViews\n2K" ]
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https://www.arxiv-vanity.com/papers/1201.4060/
[ "arXiv Vanity renders academic papers from arXiv as responsive web pages so you don’t have to squint at a PDF. Read this paper on arXiv.org.\n\n# On the efficiency and accuracy of interpolation methods for spectral codes\n\nM.A.T. van Hinsberg1    J.H.M. ten Thije Boonkkamp2    F. Toschi1 2 3    H.J.H. Clercx1 4\n22Department of Physics, Eindhoven University of Technology, PO Box 513, 5600MB Eindhoven, The Netherlands\n33Department of Mathematics and Computer Science, Eindhoven University of Technology, PO Box 513, 5600MB Eindhoven, The Netherlands\n44CNR, Istituto per le Applicazioni del Calcolo, Via dei Taurini 19, 00185 Rome, Italy\n55Department of Applied Mathematics, University of Twente, PO Box 217, 7500 AE Enschede, The Netherlands\n###### Abstract\n\nIn this paper a general theory for interpolation methods on a rectangular grid is introduced. By the use of this theory an efficient B-spline based interpolation method for spectral codes is presented. The theory links the order of the interpolation method with its spectral properties. In this way many properties like order of continuity, order of convergence and magnitude of errors can be explained. Furthermore, a fast implementation of the interpolation methods is given. We show that the B-spline based interpolation method has several advantages compared to other methods. First, the order of continuity of the interpolated field is higher than for other methods. Second, only one FFT is needed whereas e.g. Hermite interpolation needs multiple FFTs for computing the derivatives. Third, the interpolation error almost matches the one of Hermite interpolation, a property not reached by other methods investigated.\n\nI\n{AMS}\n\n65T40, 65D05\n\nnterpolation, B-spline, three-dimensional, Hermite, Fourier, spline\n\n## 1 Introduction\n\nIn recent years many studies on the dynamics of inertial particles in turbulence have focussed on the Lagrangian properties, see the review by Toschi and Bodenschatz . For particles in turbulence, but also in many other applications in fluid mechanics, interpolation methods play a crucial role as fluid velocities, rate of strain and other flow quantities are generally not available at the location of the particles, while these quantities are needed for the integration of the equations of motion of the particles.\n\nWhen a particle is small, spherical and rigid its dynamics in non-uniform flow is governed by the Maxey-Riley (MR) equation . An elaborate overview of the different terms in the MR equation and their numerical implementation can be found in the paper by Loth and a historical account was given in a review article by Michaelides . The evaluation of the hydrodynamic force exerted on the particles requires knowledge of the fluid velocity, its time derivative and gradients at the particle positions and turns out to be rather elaborate. First, the Basset history force is computationally very expensive. However, a significant reduction of cpu-time can be obtained by fitting the diffusive kernel of the Basset history force with exponential functions, as recently shown by Van Hinsberg et al. . Second, the interpolation step itself can be very time consuming and memory demanding. Especially for light particles, which have a mass density similar to the fluid density (which is, for example, sediment transport in estuaries and phytoplankton in oceans and lakes), most terms in the Maxey-Riley equation cannot be ignored and therefore also the first derivatives of the fluid velocity are needed . For this reason simulations of light particles are computationally expensive while simulations of heavy particles are less expensive. In order to achieve convergence of the statistical properties (probability distribution functions, correlation functions, multi-particle statistics, particle distributions) many particles are needed and this calls for fast and accurate interpolation methods. Therefore, our aim is to reduce the computation time for the evaluation of the trajectories of light particles substantially and make the algorithm competitive with the fast algorithms for the computation of trajectories of heavy particles in turbulence.\n\nThe incompressible Navier-Stokes equations are used to describe the turbulent flow field. In turbulence studies the Navier-Stokes equations are often solved by means of dealiased pseudo-spectral codes because of the advantage of exponential convergence of the computed flow variables. Therefore, we will focus here on interpolation methods for spectral codes.\n\nThere are many interpolation methods available . We are interested in those interpolation methods which are characterized by the following properties. First, the method must be accurate, thus we need a high order of convergence. Second, the interpolant must have a high order of continuity , with the order of continuity. Third, the method must be fast, i.e. computationally inexpensive. A very simple interpolation method is linear interpolation. This method is very fast, but unfortunately this method is relatively inaccurate and it has a low order of convergence. High order of convergence can be reached by employing Lagrange interpolation . This interpolation method has the drawback that it still has a low order of continuity for the interpolant. A solution for this problem was recently found by Lalescu et al. who proposed a new spline interpolation method. Here, the interpolant has a higher order of continuity, but the order of convergence has decreased. A method that has both a high order of convergence and a high order of continuity is Hermite interpolation . The major disadvantage of this method is that also the derivatives of the function to be interpolated are needed, these are calculated by additional Fast Fourier Transforms (FFTs) making this method computationally expensive. A remedy to this is B-spline interpolarion , which has a high order of convergence and errors comparable with the ones of Hermite interpolation. Furthermore, this method has a higher order of continuity compared with the other methods mentioned above. Normally, the transformation to the B-spline basis is an expensive step, but by making use of the spectral code it can be executed in Fourier space which makes it inexpensive. By executing this step in Fourier space the method can be optimized, resulting in smaller errors. We believe that the proposed combination of B-spline interpolation with a spectral code makes the method favorable over other traditional interpolation methods.\n\nBesides exploring the advantages of B-spline interpolation we focus on necessary conditions allowing general 3D-interpolations to be efficiently executed as successive 1D-interpolations. These conditions also carry over desired properties (like order of convergence and order of continuity) from the 1D-interpolation to the three-dimensional equivalent. Further, we provide a fast, generic algorithm to interpolate the function and its derivatives using successive 1D-interpolations.\n\nWe provide expressions for the interpolation errors in terms of the Fourier components. For this we use Fourier analysis where the interpolation method is represented as a convolution function. By doing this, the errors can be calculated as a function of the wave number. This gives insight in the behavior of interpolation, especially which components are dominant in the interpolation.\n\nThe present study may also be useful for many other applications. Examples include the computation of charged particles in a magnetic field [12, 13], but also digital filtering and applications in medical imaging [7, 14]. In the latter case interpolations are used to improve the resolution of images. Many efforts have been taken to find interpolation methods with optimum qualities . Still, it is a very active field of research. Besides the optimization of interpolation algorithms (accuracy, efficiency), the impact of different interpolation methods on physical phenomena like particle transport has been investigated in many studies [15, 16, 17].\n\nIn Section 2 we introduce the general framework and explain some one-dimen-sional interpolation methods. In Section 3 the framework is generalized to three-dimensional interpolation, and a generic algorithm is proposed for the implementation of the interpolation in Section 4. A Fourier analysis of the interpolation operator is discussed in Section 5. In Section 6 the Fourier analysis is extended to Hermite interpolation and a proof of minimal errors is given. In Section 7 our B-spline based interpolation method is introduced, and is compared with three other methods (including Hermite interpolation) in Section 8. Finally, concluding remarks are given in Section 9.\n\n## 2 Interpolation methods\n\nWe present a general framework to discuss the various interpolation methods. The goal of any interpolation method is to reconstruct the original function as closely as possible. As in many applications also some derivatives of the original function are needed, we focus on reconstructing them as well. We start with one-dimensional (1D) interpolation and subsequently, in Section 3, we generalize our framework to the three-dimensional (3D) case.\n\nLet be a 1D function that needs to be reconstructed with . In practice we only have the values of on a uniform grid, with grid spacing and knots at positions , with . After interpolation, the function is obtained which is an approximation of . Now let be the interpolation operator, so .\n\nWhen has periodic boundary conditions, it can be expressed in a Fourier series as follows\n\n u(x)=∑k∈Z^ukϕk(x),      ϕk(x)=e2πikx, (1)\n\nwith i the imaginary unit and the wave number. As the grid spacing is finite, a finite number of Fourier modes can be represented by the grid. From now on we consider to have a finite number of Fourier modes, so\n\n u(x)=kmax∑k=−kmax^ukϕk(x), (2)\n\nwhere , related to , is the maximum wave number. As we add a finite number of continuously differentiable Fourier modes we have , a property which can be used when constructing the interpolation method. In principle could be reconstructed at any point by the use of Fourier series, however in practice this would be far too time consuming and it is therefore not done, instead an interpolation is performed. is defined as the interpolant of , i.e., . We restrict ourselves to linear interpolation operators, i.e., with . This property can be used to write as\n\n ˜u(x)=kmax∑k=−kmax^uk˜ϕk(x). (3)\n\nWe focus on reconstructing by piecewise polynomial functions of degree . For each interval with we have\n\n ˜u(x)=N∑i=1aixi−1=aT¯x,      x∈(xj,xj+1),      ¯x=⎛⎜ ⎜ ⎜ ⎜ ⎜ ⎜ ⎜⎝1xx2⋮xN−1⎞⎟ ⎟ ⎟ ⎟ ⎟ ⎟ ⎟⎠. (4)\n\nHere, the vector a depends on the interval under consideration and denotes the transpose of a. The degree of the highest polynomial function for which the interpolation is still exact is denoted by . In this way we get the restriction . We consider the reconstruction of between the two neighboring grid points and . Without loss of generality we can translate and rescale so that the interval becomes the unit interval .\n\nFor Hermite interpolation the values of and of its derivatives, up to the order ( must be even), must coincide with those of at and , i.e.,\n\n dl˜udxl(0)=dludxl(0),        dl˜udxl(1)=dludxl(1),        l=0,1,..,N2−1. (5)\n\nIf the derivatives are known then . When the derivatives are not known exactly on the grid they have to be approximated by finite difference methods, as done by Lalescu et al. . Unfortunately, this method is less accurate than Hermite interpolation and .\n\nThe general framework will be illustrated with cubic Hermite interpolation for which . So the interpolation uses the function value and the first derivative in the two neighboring grid points to construct the interpolation polynomial. We have chosen this method because it is very accurate. Moreover, the second derivative, which is a piecewise linear function, gives minimal errors with respect to the -norm. This property is further discussed in Section 6.\n\nFirst, the discrete values of and possible derivatives which are given on the grid, are indicated with the vector b. In general we have\n\n b=f[u], (6)\n\nwhere the linear operator f depends on the interpolation method and maps a function onto a -vector. Second, the coefficients of the monomial basis need to be computed. Because and f are linear operators, we can write without loss of generality,\n\n aT=bTM. (7)\n\nHere, M is the matrix that defines the interpolation method. For illustration, f and M for cubic Hermite interpolation, are given by\n\n f[u]=⎛⎜ ⎜ ⎜ ⎜ ⎜ ⎜ ⎜⎝u(0)dudx(0)u(1)dudx(1)⎞⎟ ⎟ ⎟ ⎟ ⎟ ⎟ ⎟⎠,        M=⎛⎜ ⎜ ⎜⎝10−3201−21003−200−11⎞⎟ ⎟ ⎟⎠. (8)\n\nFinally, substituting relation (7) in (4) gives\n\n I[u](x)=˜u(x)=aT¯x=bTM¯x. (9)\n\nIn many applications also derivatives are needed. In order to compute the -th derivative of , the monomial basis functions should be differentiated times. In general this can be done by multiplying a by the differentiation matrix D times, so\n\n a(l)T=aTDl,          D=⎛⎜ ⎜ ⎜ ⎜ ⎜ ⎜ ⎜ ⎜ ⎜ ⎜⎝0⋯⋯⋯01⋱⋮02⋱⋮⋮⋱⋱⋱⋮0⋯0N−10⎞⎟ ⎟ ⎟ ⎟ ⎟ ⎟ ⎟ ⎟ ⎟ ⎟⎠, (10)\n\nwhere contains the coefficients for the -th derivative, obtaining\n\n dl˜udxl(x)=a(l)T¯x=bTMDl¯x=bTM(l)¯x, (11)\n\nwith . Note that the matrix D is nilpotent, since for , implying that at most derivatives can be approximated.\n\nIn conclusion, we presented a framework that is able to describe interpolation methods, which can be used to implement the interpolation methods in a straightforward way. In Section 4 it is used to generate fast algorithms for the implementation of the method.\n\n## 3 3D interpolation\n\nIn this section the 1D interpolation methods of Section 2 are extended to the 3D case. Now the scalar field depends on the vector x and a 3D interpolation needs to be carried out. Like before the interpolated field is given by and is the 3D equivalent of , so . The 3D field can be represented by a 3D Fourier series like\n\n u=∑k^ukϕk(%x), (12)\n\nwhere is given by\n\n ϕk(x)=e2πik⋅x=ϕkx(x)ϕky(y)ϕkz(z),      k=(kx,ky,kz),  x=(x,y,z), (13)\n\nand defined by (1). Again we restrict ourselves to linear interpolation operators, therefore can be written as\n\n ˜u(x)=∑k^uk˜ϕk(x), (14)\n\nwith the interpolant of , i.e., .", null, "Figure 1: Graphical description of the 3D Lagrange interpolation, using three steps of 1D interpolations for the case N=4. First, N2 1D interpolations are carried out in the x-direction (crosses). Second, N interpolations are carried out in y-direction (dots in the right figure) and from these N results finally one interpolated value is derived in z-direction (triangle).\n\nThe 3D interpolation for a scalar field is carried out applying three times 1D interpolations, see Fig. 1. The interpolation consists of three steps, in which the three spatial directions are interpolated one after each other. The order in which the spatial directions are interpolated does not matter. Building the 3D interpolation out of 1D interpolations saves computing time. It can be done for all interpolation methods as long as the following two conditions are met. First, the operator must be linear. Second, the following condition must be satisfied\n\n ˜ϕk≡I3[ϕk]=I3[ϕkxϕkyϕkz]=I[ϕkx]I[ϕky]I[ϕkz]=˜ϕkx˜ϕky˜ϕkz, (15)\n\nwhich is the case for almost all interpolation methods. Property (15) can be used to prove that properties of the 1D operator carry over to the 3D operator , for example, the order of convergence and the order of continuity.\n\nNext, relations (9) and (11) are extended to the 3D case. Like before, we start with storing some values of (given by the spectral code) and possible derivatives in the third-order tensor B. In the same fashion as relation (6) one gets\n\n B=fz[fy[fx[u]]], (16)\n\nwhere one element of tensor B is defined like\n\n Bi1i2i3=fz[fy[fx[u]i1]i2]i3. (17)\n\n, and are similar to operator f but now working in a specified direction. For Hermite interpolation they are given by\n\n fx[u]=⎛⎜ ⎜ ⎜ ⎜ ⎜⎝u(0,y,z)∂u∂x(0,y,z)u(1,y,z)∂u∂x(1,y,z)⎞⎟ ⎟ ⎟ ⎟ ⎟⎠,    fy[u]=⎛⎜ ⎜ ⎜ ⎜ ⎜ ⎜⎝u(x,0,z)∂u∂y(x,0,z)u(x,1,z)∂u∂y(x,1,z)⎞⎟ ⎟ ⎟ ⎟ ⎟ ⎟⎠,    fz[u]=⎛⎜ ⎜ ⎜ ⎜ ⎜⎝u(x,y,0)∂u∂z(x,y,0)u(x,y,1)∂u∂z(x,y,1)⎞⎟ ⎟ ⎟ ⎟ ⎟⎠. (18)\n\nThe interpolation is carried out in a similar way as sketched in Fig. 1. Similarly to (9), can be represented as\n\n I3[u](x)=˜u(x)=% B¯×1(M¯x)¯×2(M¯y)¯×3(M¯z), (19)\n\nwhere M is still the matrix for 1D interpolation, and are defined like which is given by relation (4). Further, denotes the -mode vector product , like\n\n A=B¯×nf,      Ai1⋯in−1in+1⋯iN=∑inBi1⋯iNfin, (20)\n\nwhere denotes the order of tensor B. In this way tensor A is one order less than tensor B. Because we employ three of these -mode vector products the third-order tensor B reduces to a scalar. Furthermore, each of these -mode vector products corresponds to an interpolation in one direction, see also Fig. 1. For a general derivative one gets\n\n ∂i+j+k˜u∂xi∂yj∂zk(x)=B¯×1(M(i)¯x)¯×2(M(j)¯y)¯×3(M(k)¯z). (21)\n\nNote that the matrix M does not necessarily have to be the same for the different directions , and . One could choose different interpolation methods when for example Chebyshev polynomials are used in one direction. In this case the grid is nonuniform in this direction and therefore not all interpolation methods can be used.\n\nFinally, when the scalar field becomes a vector field , the three components of u can be interpolated separately. This can be written in short by a fourth order tensor B where the last dimension contains the three components of u. In this way the equations for the new tensor B remain the same as given above.\n\n## 4 Implementation\n\nRelations (19) and (21) provide a good starting point for an efficient implementation of the interpolation. We focus on interpolating a 3D vector field and on calculating all its first derivatives (which are needed in many applications like the computation of the trajectories of inertial particles). The matrices M and only need to be computed once, which can be done prior to interpolation. Second, the vectors , and have to be computed which only needs to be done once for each position of interpolation. In Table 1 we keep track of all the computed quantities. Here, the computational costs for evaluating all the components is shown where one flop denotes one multiplication with one addition. We show the number of flops for the general case and for . The main idea is to reduce the order of the tensors as soon as possible in order to generate an efficient method.\n\nIn order to determine how efficient the algorithm is, one can compare the computational costs against a lower bound. The lower bound we use is related to the size of B which is for a vector field u. In order to be able to use all the information in tensor B, flops are needed. For large one finds that the algorithm of Table 1 is only a factor 2 less efficient than this lower bound.\n\nWe also compare our algorithm with the one proposed by Lekien and Marsden , which uses Hermite interpolation with . Our algorithm has less restrictions and shows a slightly better computational performance (for ). The algorithm of Lekien and Marsden consists of two steps. First, they calculate the coefficients for the polynomial basis. Second, the values at the desired location are calculated. They claim that their method is beneficial when the derivatives are needed or the interpolation needs to be done multiple times for the same interval, because only the second step needs to be executed multiple times. Our method does not have the first step, therefore it has no restrictions, nevertheless the computation of the values and the first derivatives is slightly faster than for Lekien and Marsden, even when considering only the second step. The total costs of their second step is bounded by flops (4 times flops, for the computation of the values and the first derivatives). From Table 1 we can conclude that our method needs less flops for the same computations.\n\n## 5 Fourier analysis\n\nIn this section the interpolation operator is expressed in terms of a convolution. In this way properties of the interpolation method like the order of continuity of the interpolated field and the magnitude of the errors can be shown in the Fourier domain. We start with the interpolation of 1D functions and subsequently, it can be extended to the 3D case.", null, "Figure 2: Sketch of linear interpolation as a convolution. The pins represent delta functions with the height equal to its prefactor. On the left side is a visualization in real space and on the right side in Fourier space.\n\nBefore we start with the derivation, we rescale the variable by dividing it by , so that the new grid spacing equals unity. From now on we work with the rescaled grid where and , so for . Furthermore we introduce the dimensionless wave number and is similarly defined as , see (1). For Hermite interpolation the derivation is somewhat more complex because also the derivatives are used and therefore it is postponed to Section 6. We focus on interpolation methods where contains the values of at the nearest grid points of with local ordering. Thus and is given by\n\n xj = ⌊x−N2+j⌋,            j=1,2,⋯,N, (22)\n\nwhere denotes the nearest lower integer. The interpolation methods can be described by the matrix M, with elements , see relation (9). This relation can also be written as\n\n ˜u(x) = N∑j=1Cj(x−xj)u(xj), (23)\n\nwith defined by (22) and where is given by\n\n Cj(x+N2−j) = {∑Ni=1Mj,ixi−1% for 0≤x<1,0elsewhere. (24)\n\nRelation (23) can be rewritten by using the sifting property of the delta function, like\n\n ˜u(x) = N∑j=1Cj(x−xj)∫∞−∞u(y)δ(y−xj)dy. (25)\n\nThis can be further reformulated by subtracting the argument of the delta function from the argument of , as\n\n ˜u(x) = ∫∞−∞u(y)N∑j=1δ(y−xj)Cj(x−y)dy (26) = ∫∞−∞u(y)D(y)C(x−y)dy,\n\nwith and given by\n\n C(x)=N∑j=1Cj(x),          D(x)=∑i∈Zδ(x−i). (27)\n\nIn relation (26) the delta functions can be replaced by the function , which is a train of delta functions because the functions only have a support of length unity, see (24). Finally, the interpolation can be written like\n\n ˜u=(uD)∗C, (28)\n\nwith denoting the convolution product. Here, the convolution function depends on the interpolation matrix M, see Fig. 2.\n\nAs a consequence of relation (28), if the function is continuous up to the -th derivative then is also continuous up to the -th derivative. Even stronger, the order of continuity of the function is equal to the order of continuity of . Furthermore, by the use of relation (28) exact interpolation can be constructed as well 111Exact interpolation can be accomplished by for and zero elsewhere. In this way only the original Fourier component is filtered out of the spectrum. Note that in this case has infinite support..\n\nIn the following of this section we will discuss the interpolation error. Before proceeding we need to proof the following theorem.\n\nTheorem. for . Here is the error in mode , and is the inner product related to the -norm defined by\n\n ⟨f,g⟩2=∫m0f(x)g∗(x)dx,      ∥f∥22=⟨f,f⟩2=∫m0f(x)f∗(x)dx. (29)\n\nThe asterisk denotes complex conjugation.\n\n ˜ϕκ=I[ϕκ]=(ϕκD)∗C. (30)\n\nSecond, we take the Fourier transform of , for some fixed , i.e.,\n\n F[˜ϕκ0](k) = F[(ϕκ0D)∗C](k)=(F[ϕκ0]∗F[D])(k)F[C](k) (31) = m∑i∈Zδ(k−(i+κ0))F[C](i+κ0),\n\nwhich results in a train of delta functions with the perfector given by , see Fig. 2, and denotes the Fourier transform given by\n\n F[g](k):=∫∞−∞g(x)e−2πikxdx. (32)\n\nFor linear interpolation these functions are shown in Fig. 2. Next, can be written as . Trivially for . Furthermore, consists of a discrete set of Fourier components, see relation (31). Using this relation, one can show that no common Fourier components exist for and or for . Therefore , and for implying as claimed.\n\nCorollary. The orthogonality is important to estimate errors. When the error in is computed as , it can be rewritten like , which allows easy and straightforward computation of the errors.\n\nNext, the error in one Fourier component is calculated. In this derivation we make use of the fact that can be written by a sum of Fourier components, see Fig. 2 and relation (31). The relative error in one Fourier component can be written as\n\n ∥∥˜ϕκ−ϕκ∥∥22∥ϕκ∥22=∥eκ∥22m = 1m∥∥ ∥∥−ϕκ+∑i∈ZF[C](κ+i)ϕκ+i∥∥ ∥∥22 (33) = (F[C](κ)−1)2+∑i≠0(F[C](κ+i))2.\n\nFrom this expression one can see that the error can be computed directly from . The same can be done for the error in the -th derivative; . The idea is to take the derivatives of the individual Fourier components which results in\n\n ∥∥e(l)κ∥∥22∥∥ϕ(l)κ∥∥22=(F[C](κ)−1)2+∑i≠0(κ+iκ)2l(F[C](κ+i))2. (34)\n\nThe extension to the 3D case is rather straightforward and is therefore not reported here. The basic idea is to create 3D functions by multiplying the 1D components, this can be done for all functions and the basic equations remain the same.\n\n## 6 Hermite interpolation\n\nIn this section we extend the theory of Section 5 to Hermite interpolation. We also show some special properties that hold for Hermite interpolations. Especially, we examine the case . For this case the second derivative becomes a piecewise linear function. Comparison with the actual second derivative shows that this piecewise linear function is optimal with respect to the -norm.\n\nIn order to extend the theory of Section 5 to Hermite interpolation with even the same procedure is followed as in Section 5. Analogous to (23), can be written as\n\n ˜u(x) = 1∑j=0N/2∑l=1Cj,l(x−xj)dl−1udxl−1(xj), (35)\n\nwhere and are given by\n\n Cj,l(x−j) = {∑Ni=1Ml+jN2,ixi−1for 0≤x<10elsewhere,        l∈1,2,⋯,N2, xj = ⌊x⌋+j,                         j∈0,1. (36)\n\nAgain following similar steps as in Section 5, can be rewritten as\n\n ˜u(x) = 1∑j=0N/2∑l=1Cj,l(x−xj)∫∞−∞dl−1udxl−1(y)δ(y−xj)dy (37) = N/2∑l=1∫∞−∞dl−1udxl−1(y)D(y)Cl(x−y)dy,\n\nwhere is given by relation (27) and is given by . In short, can be written as\n\n ˜u=N/2∑l=1(dl−1udxl−1D)∗Cl, (38)\n\nsimilar to relation (28). Here one can see that for Hermite interpolation multiple convolution functions are needed which correspond to the derivatives and the function itself. Replacing by in (38) gives\n\n ˜ϕκ=I[ϕκ]=(ϕκD)∗N/2∑l=1(2πiκ)l−1Cl. (39)\n\nIn this way we find a similar expression as relation (30), where has to be replaced by . In conclusion, relation (33) and (34) can still be used.\n\nProperty. For the error in the first derivative we have the following property:\n\n ⟨e(1),1⟩2=0, (40)\n\nwhere the inner product is defined on the unit interval, i.e.,\n\n ⟨f,g⟩2=∫10f(x)g∗(x)%dx. (41)\n\nFurthermore the error in the -th derivative, , is given by\n\n e(l)=dl˜udxl(x)−dludxl(x). (42)\n\nProof. One can rewrite, part of the interpolation conditions for Hermite interpolation (5) in the following way\n\n ˜u(1)−˜u(0)=u(1)−u(0)  ⇔  ∫10d˜udxdx=∫10dudxdx. (43)\n\nHere two interpolation conditions give one new condition which is equivalent to relation (40).\n\nCorollary. Property (40) shows that the error in the first derivative does not have a constant component. Therefore the constant component is exact with respect to the -norm\n\nProperty. For the error in the second derivative in case of we have\n\n (44)\n\nProof. One can rewrite the interpolation conditions (5) for in the following way\n\n ˜u(1)−˜u(0)−˜u′(0)=u(1)−u(0)−u′(0) ⇔ ∫10∫α0d2˜udx2dxdα=∫10∫α0d2udx2dxdα, ˜u′(1)−˜u′(0)=u′(1)−u′(0) ⇔ ∫10d2˜udx2dx=∫10d2udx2dx. (45)\n\nThe first relation in (44) follows immediately from the second condition in (45). The second relation in (44) is derived in the following way\n\n ∫10∫α0d2˜udx2dxdα=∫10∫α0d2udx2dxdα, α∫α0(d2˜udx2(x)−d2udx2(x))dx∣∣ ∣∣α=1α=0−∫10(d2˜udx2(α)−d2udx2(α))αdα=0, ∫10(d2˜udx2(x)−d2udx2(x))xdx=0. (46)\n\nHere, the first step is integration by parts and the second step uses the second relation of equation (45).\n\nCorollary. Relation (45) implies that does not have a constant component, neither a linear component. For the second derivative is a linear function, and this means that there is no better approximation in the -norm of this second derivative with a piecewise linear function. This makes Hermite interpolation very interesting as a reference case, because we now have proven that the error is minimal for this case.\n\n## 7 B-spline interpolation\n\nIn this section we start with explaining B-spline interpolation. The idea is to create an as smooth as possible interpolant. Later it is shown how the pseudo-spectral code can be used to efficiently execute this interpolation method. Furthermore, the interpolation method is optimized to create small errors in the -norm. We start with giving the B-spline convolution functions after which their matrix representation is given and finally the transformation to the B-spline basis functions is derived.\n\nIn a spectral code FFTs are applied to transform data from real space to Fourier space and backwards. These FFTs are the most expensive step in the simulation and therefore we want to keep the number of FFTs needed minimal. This is the reason why Hermite interpolation is not a good option, since extra FFTs are needed for the computation of the derivatives. An alternative is B-spline interpolation.\n\nWe require high order of continuity of the interpolant. The highest order of continuity that can be obtained for the interpolant with piecewise polynomial functions of degree is . In this way the interpolant still matches the original function at the grid points . Moreover, one can immediately see that , where is the highest degree of a polynomial test function for which the interpolation is still exact. This high level of continuity can be achieved by using B-spline functions . The first four uniform B-spline basis functions are shown in Fig. 3. These functions can be generated by means of convolutions in the following way\n\n B(1)(x) = {1for    −0.5≤x<0.5,0elsewhere, B(2) = B(1)∗B(1), B(3) = B(2)∗B(1), ⋮ B(N) = B(N−1)∗B(1). (47)\n\nThese functions have the property that the -th function is of degree and is . Furthermore, the B-spline basis functions have local support of length . The B-spline functions can be seen as convolution functions introduced in Section 5 and have a matrix representation. The relation between the functions and the matrix representation is similar to relation (24) and (27), and is given by\n\n B(N)(x) = N∑j=1B(N),j(x), B(N),j(x+N2−j) = {∑Ni=1M(N),i,jxi−1for 0≤x<1,0elsewhere. (48)\n\nThe matrix representation for the first four B-spline functions is as follows \n\n M(1) = (1), M(2) = (1−101), M(3) = 12!⎛⎜⎝1−2112−2001⎞⎟⎠, M(4) = 13!⎛⎜ ⎜ ⎜⎝1−33−140−63133−30001⎞⎟ ⎟ ⎟⎠. (49)\n\nIn general we have \n\n M(N),i,j=1(N−1)!QN−jN−1N∑s=i(−1)s−iQs−iN(N−s)N−j,       i,j=1,2,⋯,N, (50)\n\nwith given by\n\n Qin=n!i!(n−i)!=(ni). (51)\n\nWe still need to express in terms of B-spline basis functions and thus find the transform from real space to the B-spline basis. Because the inverse transform from the B-spline basis to real space is somewhat easier, we start with this transformation first. From now on we omit the subindex . The coefficients of the B-spline basis are called , and can be derived from it by the discrete convolution in the following way, . Here, is given by\n\n BD(x)={B(x)for x\n\nand the discrete convolution is given by\n\n (g∗Dh)(x)=m∑y=0g(y)h" ]
[ null, "https://media.arxiv-vanity.com/render-output/3134487/x1.png", null, "https://media.arxiv-vanity.com/render-output/3134487/x2.png", null ]
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https://skoolers.com/physics/
[ "Physics\n\nPhysics Syllabus\n\nCXC / CSEC PAST PAPERS\n\nVideoes\n\nSECTION A: MECHANICS\n\nScientific Method\n\nGalileo\n\nSimple Pendulum\n\nMeasurement\n\n• Express the result of a measurement or a calculation to an appropriate number of significant figures;\n• Discuss the possible types and sources of error in any measurement;\n• Use a variety of instruments to measure different quantities;\n• Assess the suitability of instruments on the basis of sensitivity, accuracy and range;\n• Apply the formula for density p=m/v.\n\nVectors\n\nStatics\n\nForces, F\n\n• Explain the effects of forces;\n• Identify types of forces;\n• Determine the weight of objects;\n• Show how derived quantities and their related units are produced;\n• Recall the special names given to the units for some derived quantities;\n• Express derived units using the index notation;\n• Identify situations in which the application of a force will result in a turning effect.\n\nTurning Forces\n\n• Define the moment of a force T;\n• Apply the principle of moments;\n• Explain the action of common tools and devices as levers;\n• Determine the location of the centre of gravity of a body;\n• Relate the stability of an object to the position of its centre of gravity and its weight.\n\nDeformation\n\n•  Investigate the relationship between extension and force;\n• Solve problems using Hooke’s law.\n\nDynamics: Motion in a Straight Line\n\n• Define the terms: distance, displacement, speed, velocity, acceleration;\n• Apply displacement-time and velocity-time graphs.\n\nAristotle\n\n• Discuss Aristotle’s arguments in support of his “law of motion”, that is vαF.\n\nNewton’s Laws\n\n• State Newton’s three laws of motion;\n• Use Newton’w laws to explain dynamic systems;\n• Define linear momentum;\n• Describe situations that demonstrate the law of conservation of linear momentum;\n• Apply the law of conservation of linear momentum.\n\nEnergy\n\nForms of energy\n\n• Define Energy;\n• Identify the various forms of energy;\n• Describe the energy transformations in a given situation;\n• Apply the relationship: work = force x displacement;\n• Discuss the use of energy from alternative sources, and its importance to the Caribbean.\n\nPotential Energy, Ep\n\n• Define potential energy;\n• Calculate the change in graviational potential energy using: ΔEp = mgΔh.\n\nKinetic Energy, Ek\n\n• Define kinetic energy;\n• Calculate kinetic energies using the expression: Ek = ½ mv2.\n\nConservation\n\n• Apply the law of conservation of energy.\n\nPower, P\n\n• Define power and apply definition;\n• Explain the term efficiency;\n• Calculate efficiency in given situations.\n\nHydrostatics\n\n• Define pressure and apply definition;\n• Relate the pressure at a point in a fluid to its depth and the density;\n• Apply Archimedes’ principle to predict whether a body would float or sink in a given fluid.\n\nSECTION B: THERMAL PHYSICS AND KINETIC THEORY\n\nNature of Heat\n\n• Differentiate between caloric and kinetic theories of heat as they existed in the eighteenth century;\n• Discuss the role of Joule’s experiments in establishing the principle of conservation of energy.\n\nMacroscopic Properties and Phenomena\n\nTemperature, T\n\n• Relate temperature to the direction of net thermal energy transfer;\n• Identify physical properties which may vary with temperature and may be used as a basis for measuring temperature;\n• Relate the use of a thermometer to its design;\n• Define the fixed points on the Celsius scale;\n• Relate the temperature of a body to the kinetic energy of molecules.\n\nPhases of Matter\n\n• Distinguish among solids, liquids and gases;\n• Use the kinetic theory to explain the different macroscopic properties of solids, liquids and gases.\n\nExpansion\n\n• Explain observations of the effects of thermal expansion.\n\nGas Laws\n\n• Relate graphs of pressure or volume against temperature to establishment of the Kelvin temperature scale;\n• Use the relationship between Kelvin and Celsius scale. T/K = 0 degrees C+ 273;\n• Apply the gas laws;\n• Give the qualitative explanations of the gas laws in terms of kinetic theory.\n\nThermal Measurements\n\nSpecific Heat Capacity, c\n\n• Distinguish between specific heat capacity, ‘c’ and heat capacity ‘C’;\n• Apply the relationship EH = mcθ or EH = mcΔT;\n• Determine the specific heat capacity of metals and liquids.\n\nSpecific Latent Heat, l\n\n• Demonstrate that temperature remains constant during a phase change;\n• Apply the relationship EH = ml;\n• Determine the specific latent heat of vaporization lv, and fusion, lf of water;\n• Distinguish between evaporation and boiling.\n\nTransfer of Thermal Energy\n\n• Explain the transfer of thermal energy by conduction;\n• Explain the transfer of thermal energy by convection;\n• Explain the transfer of thermal energy by radiation;\n• Conduct experiments to investigate the factors on which absorption and emission of radiation depend;\n• Recall that good absorbers are good emitters;\n• Relate the principles of thermal energy transfer to the design of other devices.\n\nSECTION C: WAVES AND OPTICS\n\nWave Motion\n\nTypes of Waves\n\n• Differentiate between types of waves.\n\nWave Parameters\n\n• Apply speed, frequency, wavelength, period and amplitude;\n• Represent transverse and longitudinal waves in displacement position and displacement-time graphs.\n\nSound\n\nProduction and Propagation\n\n• Describe how sound is produced and propagated in a medium;\n• Relate the terms ‘pitch’ and ‘loudness’ to wave parameters.\n\nSpeed of Sound\n\n• Apply the speed of sound to practical situations;\n• Cite evidence that sound waves reflect, refract, diffract and interfere;\n• Describe the use of ultrasound.\n\nElectromagnetic Waves\n\n• State the properties of e.m. waves;\n• Differentiate between types of e.m.waves in terms of their wavelengths;\n• Identify a source and use of each type of e.m.wave.\n\nLight Waves\n\nWave Particle Duality\n\n• Compare the rival theories of light held by scientists;\n• Conduct a Young’s double slit experiment to show that light is a wave.\n\nRays of Light\n\n• Explain why diffraction of light is not normally observed;\n• Apply the principle that light travels in straight lines.\n\nReflection\n\n• Apply the laws of reflection;\n• Describe the formation of images in a plane mirror.\n\nRefraction\n\n• Give examples of observations which indicate that light can be refracted;\n• Describe the refraction of light rays;\n• Describe how a prism may be used to produce a spectrum;\n• Apply Snell’s Law.\n\nCritical Angle and Total Internal Reflection\n\n• Explain ‘critical angle’ and ‘total internal reflection’;\n• Relate critical angles to total internal reflection;\n• Draw diagrams illustrating applications of total internal reflection.\n\nLenses\n\nAction of Lenses\n\n• Illustrate the effect of converging and diverging lenses on a beam of parallel rays;\n• Define the terms: principal axis, principal focus, focal length, focal plane, magnification.\n\nImage Formation\n\n• Differentiate between real and virtual images;\n• Apply the equations for magnification;\n• Determine the focal length of a converging lens.\n\nSECTION D: ELECTRICITY AND MAGNETISM\n\nElectrostatics\n\nElectric Charge, Q\n\n• Explain the charging of objects;\n• Describe the forces that electric charges exert on each other;\n• Explain charging by induction.\n\nElectric Fields\n\n• Define an electric field;\n• Describe one hazard and one useful application of static charge.\n\nCurrent Electricity\n\n• Distinguish between conductors and insulators;\n• State than an electric current in a metal consists of a flow of electrons;\n• Differentiate between electron flow and conventional current;\n• State the unit of electrical current;\n• apply the relationship Q= It.\n\nAlternating Current\n\n• Differentiate between direct and alternating currents;\n• Analyse current-time or voltage time graphs.\n\nElectrical Quantities\n\nPower, P and Energy, E\n\n• Cite examples of conversion of electrical energy to other forms and vice versa;\n• Apply the relationship V= E/Q;\n• Apply the relationship P= IV;\n• Discuss the importance of conserving electrical energy and the means of doing so.\n\nCircuit and Components\n\nCircuit Diagrams\n\n• Use symbols to construct circuit diagrams;\n• Differentiate between series and parallel circuits.\n\nCells\n\n• Explain the functions of the various parts of a zinc-carbon cell;\n• Distinguish between primary and secondary cells;\n• Draw a circuit diagram to show how a secondary cell can be recharged.\n\nI-V Relationships\n\n• Investigate the relationship between current and potential difference.\n\nResistance, R\n\n• Explain the concept of resistance;\n• Apply the relationship R= V/I;\n• Explain why it is necessary for an ammeter to have a very low resistance;\n• Explain why it is necessary for a voltmeter to have very high resistance;\n• Solve problems involving series and parallel resistance;\n• Solve problems involving series, parallel and series parallel circuits.\n\nElectricity in the Home\n\n• Discuss the reasons for using parallel connections of domestic appliances;\n• Explain the purpose of a fuse or circuit breaker and the earth wire;\n• Select a fuse or circuit breaker of a suitable current rating for a given appliance;\n• State the adverse effects of connecting electrical appliances to incorrect or fluctuating voltage supplies.\n\nElectronics\n\n• Describe how a semi-conductor dioxide can be used in half wave rectification;\n• Differentiate between direct current from batteries and rectified alternating current by a consideration of V-t graphs for both cases.\n\nLogic Gates\n\n• Recall the symbols for AND, OR, NOT, NAND, NOR logic gates;\n• State the function of truth gates with the aid of truth tables;\n• Analyse the circuits involving the combinations of not more than three logic gates;\n• Discuss the impact of electronic and technological advances on society.\n\nMagnetism\n\nPermanent Magnets\n\n• Differentiate between magnetic and non-magnetic materials;\n• Explain how a magnet can attract an unmagnetised object;\n• Distinguish between materials used to make “permanent” and “temporary” magnets;\n• Identify the poles of a magnetic dipole.\n\nMagnetic Forces\n\n• Investigate the forces between magnetic poles;\n• Define a magnetic field;\n• Map magnetic fields.\n\nElectromagnetism\n\n• Conduct simple experiments to investigate the magnetic field pattern around current carrying conductors;\n• Apply suitable rules which relate the direction of current flow to the direction of the magnetic field;\n• Describe a commercial application of an electromagnet.\n\nElectromagnetic Force\n\n• Conduct an experiment which demonstrates the -existence of a force on a current-carrying conductor placed in a magnetic field;\n• Sketch the resultant magnetic flux pattern when a current carrying wire is placed perpendicular to a uniform magnetic field;\n• Apply Flemming’s left hand motor rule;\n• Identify factors that affect the force on a current-carrying conductor in a magnetic field.\n\nMotors\n\n• Explain the action of a D.C motor.\n\nInduced e.m.f\n\n• Describe simple activities which demonstrate an induced e.m.f;\n• Conduct simple experiments to show the magnitude of an induced e.m.f;\n• Predict the direction of induced current given the direction of motion of the conductor and that of the magnetic field;\n• Explain the action of the A.C. generator.\n\nTransformers\n\n• Explain the principle of transformation of a transformer;\n• State the advantages of using a.c. for transferring electrical energy;\n• Apply the ideal transformer formula Pout = Pin\n\nSECTION E: THE PHYSICS OF THE ATOM\n\nModels of the Atom\n\n• Describe the work done in establishing the modern view of the atom;\n• Describe the Geiger-Marsden experiment.\n\nStructure of the Atom\n\nParticles in the Atom\n\n• Sketch the structure of simple atoms;\n• Compare the mass and charge of the electron with the mass and charge of the proton;\n• Explain why the atom is normally neutrally stable;\n• Apply the relationship: A = Z+N;\n• Explain what is meant by the term “isotope”;\n• Relate the shell model of the atom to the periodic table.\n\n• Describe Marie Curie’s work in the field of radioactivity;\n• State the nature of the three types of radioactive emissions;\n• Describe experiments to compare the ranges of α, β and γ emission;\n• Describe the appearance of the tracks of radioactive emissions in a cloud chamber;\n• Predict the effects of magnetic and electric fields on the motion of α and β particles and γ rays;\n• Interpret nuclear reactions in the standard form;\n• Conduct an activity to demonstrate the nature of radioactive decay;\n• Recall that the decay process is independent of the conditions external to the nucleus.\n\nHalf-life\n\n• Use graphs of random decay to show that such processes have constant half-lives;\n• Solve problems involving half-life." ]
[ null ]
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https://www.komal.hu/feladat?a=feladat&f=C981&l=en
[ "", null, "Mathematical and Physical Journal\nfor High Schools\nIssued by the MATFUND Foundation\n Already signed up? New to KöMaL?\n\n#", null, "Problem C. 981. (March 2009)\n\nC. 981. The captain of a ship wrote down the formula", null, "in his log book for the distance to the horizon. Unfortunately, the number p got blurred. d denotes the distance to the horizon in kilometers, and h stands for the height of the eyes of the observer above sea level in metres. Determine a value of p that provides a practicable formula. (Use 6370 km for the radius of the Earth.)\n\n(5 pont)\n\nDeadline expired on April 15, 2009.\n\nSorry, the solution is available only in Hungarian. Google translation\n\nMegoldás. Készítsünk ábrát.", null, "Pitagorasz-tétellel:", null, ", amiből", null, ". Ha h kicsi és R=6370 km, akkor", null, "km.\n\n### Statistics:\n\n 171 students sent a solution. 5 points: 79 students. 4 points: 18 students. 3 points: 55 students. 2 points: 11 students. 1 point: 1 student. 0 point: 2 students. Unfair, not evaluated: 5 solutionss.\n\nProblems in Mathematics of KöMaL, March 2009" ]
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https://socratic.org/questions/what-is-the-derivative-of-f-x-sin-1-x
[ "# What is the derivative of f(x)=sin^-1(x) ?\n\nSep 11, 2014\n\nMost people remember this\n$f ' \\left(x\\right) = \\frac{1}{\\sqrt{1 - {x}^{2}}}$\nas one of derivative formulas; however, you can derive it by implicit differentiation.\n\nLet us derive the derivative.\nLet $y = {\\sin}^{- 1} x$.\n\nBy rewriting in terms of sine,\n$\\sin y = x$\n\nBy implicitly differentiating with respect to $x$,\n$\\cos y \\cdot \\frac{\\mathrm{dy}}{\\mathrm{dx}} = 1$\n\nBy dividing by $\\cos y$,\n$\\frac{\\mathrm{dy}}{\\mathrm{dx}} = \\frac{1}{\\cos} y$\n\nBy $\\cos y = \\sqrt{1 - {\\sin}^{2} y}$,\n$\\frac{\\mathrm{dy}}{\\mathrm{dx}} = \\frac{1}{\\sqrt{1 - {\\sin}^{2} y}}$\n\nBy $\\sin y = x$,\n$\\frac{\\mathrm{dy}}{\\mathrm{dx}} = \\frac{1}{\\sqrt{1 - {x}^{2}}}$" ]
[ null ]
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