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http://interactiverhythm.com/f8v3zz9e/431ff7-eigenvectors-of-orthogonal-matrix-are-orthogonal
[ "Some Elements Are Metal Like Iron Gold And Silver, Plus Size Clothing Melbourne, National Bird Of Albania, Michael Bernstein Vnsny, Dendrobium Kingianum Watering, Jefferson County Alabama Board Of Education Phone Number, A Thousand Kisses Deep Lyrics Meaning, Qtile Vs Bspwm, \" />", null, "", null, "# eigenvectors of orthogonal matrix are orthogonal\n\nWednesday, December 9th, 2020\n\nThis website’s goal is to encourage people to enjoy Mathematics! Range, Null Space, Rank, and Nullity of a Linear Transformation from $\\R^2$ to $\\R^3$, How to Find a Basis for the Nullspace, Row Space, and Range of a Matrix, The Intersection of Two Subspaces is also a Subspace, Rank of the Product of Matrices $AB$ is Less than or Equal to the Rank of $A$, Find a Basis and the Dimension of the Subspace of the 4-Dimensional Vector Space, Show the Subset of the Vector Space of Polynomials is a Subspace and Find its Basis, Find a Basis for the Subspace spanned by Five Vectors, Prove a Group is Abelian if $(ab)^2=a^2b^2$, Dimension of Null Spaces of Similar Matrices are the Same. Orthogonal Eigenvectors Suppose P1, P2 € R2 are linearly independent right eigenvectors of A E R2x2 with eigenvalues 11, 12 E R such that 11 # 12. MATH 340: EIGENVECTORS, SYMMETRIC MATRICES, AND ORTHOGONALIZATION 5 By our induction hypothesis, there exists an orthogonal matrix Q such that QtBQ is diagonal. I have a Hermitian matrix, and I would like to get a list of orthogonal eigenvectors and corresponding eigenvalues. Eigenvalues and Eigenvectors The eigenvalues and eigenvectors of a matrix play an important part in multivariate analysis. All identity matrices are an orthogonal matrix. Polynomial $x^4-2x-1$ is Irreducible Over the Field of Rational Numbers $\\Q$. So the determinant of an orthogonal matrix must be either plus or minus one. Christa. Eigen decompositions tells that $U$ is a matrix composed of columns which are eigenvectors of $A$. This is an elementary (yet important) fact in matrix analysis. Eigenvectors of Symmetric Matrices Are Orthogonal - YouTube Course Hero is not sponsored or endorsed by any college or university. Again, as in the discussion of determinants, computer routines to compute these are widely, available and one can also compute these for analytical matrices by the use of a computer algebra, This discussion applies to the case of correlation matrices and covariance matrices that (1), have more subjects than variables, (2) have variances > 0.0, and (3) are calculated from data having. But we have 2 special types of matrices Symmetric matrices and Hermitian matrices. This site uses Akismet to reduce spam. Matrices of eigenvectors (discussed below) are orthogonal matrices. Inner Product, Norm, and Orthogonal Vectors. Let $A=\\begin{bmatrix} 1 & -1\\\\ 2& 3 \\end{bmatrix}.$ ... For approximate numerical matrices m, the eigenvectors are normalized. . Your email address will not be published. Finding Eigenvalues and Eigenvectors : 2 x 2 Matrix Example - Duration: 13:41. patrickJMT 1,472,884 views. I know that Matlab can guarantee the eigenvectors of a real symmetric matrix are orthogonal. ... Constructing an Orthogonal Matrix from Eigenvalues - Duration: 10:09. Proof. In fact, for a general normal matrix which has degenerate eigenvalues, we can always find a set of orthogonal eigenvectors as well. PCA of a multivariate Gaussian distribution centered at (1,3) with a standard deviation of 3 in roughly the (0.866, 0.5) direction and of 1 in the orthogonal direction. The extent of the stretching of the line (or contracting) is the eigenvalue. In fact, PTP == 2 4 122 −2−12 2−21 3 5 2 4 1−22 2−1−2 22 1 3 5= 2 4 900 090 009 3 5: . can be mathematically decomposed into a product: characteristic vectors or latent vectors. ) Given the eigenvector of an orthogonal matrix, x, it follows that the product of the transpose of x and x is zero. Symmetric matrices have n perpendicular eigenvectors and n real eigenvalues. The orthogonal matrix has all real elements in it. The above matrix is skew-symmetric. 49:10. When I use [U E] = eig(A), to find the eigenvectors of the matrix. Now without calculations (though for a 2x2 matrix these are simple indeed), this A matrix is . Then show that the nullity of $A$ is equal to... Is a Set of All Nilpotent Matrix a Vector Space? How can I demonstrate that these eigenvectors are orthogonal to each other? Condition that Vectors are Linearly Dependent/ Orthogonal Vectors are Linearly Independent, Determine the Values of $a$ such that the 2 by 2 Matrix is Diagonalizable, Sequence Converges to the Largest Eigenvalue of a Matrix, Eigenvalues of Real Skew-Symmetric Matrix are Zero or Purely Imaginary and the Rank is Even, Properties of Nonsingular and Singular Matrices, Symmetric Matrices and the Product of Two Matrices, Find Values of $h$ so that the Given Vectors are Linearly Independent, Linear Combination and Linear Independence, Bases and Dimension of Subspaces in $\\R^n$, Linear Transformation from $\\R^n$ to $\\R^m$, Linear Transformation Between Vector Spaces, Introduction to Eigenvalues and Eigenvectors, Eigenvalues and Eigenvectors of Linear Transformations, How to Prove Markov’s Inequality and Chebyshev’s Inequality, How to Use the Z-table to Compute Probabilities of Non-Standard Normal Distributions, Expected Value and Variance of Exponential Random Variable, Condition that a Function Be a Probability Density Function, Conditional Probability When the Sum of Two Geometric Random Variables Are Known, Determine Whether Each Set is a Basis for $\\R^3$. Suppose that $n\\times n$ matrices $A$ and $B$ are similar. But as I tried, Matlab usually just give me eigenvectors and they are not necessarily orthogonal. 5 years ago. Let y be eigenvector of that matrix. To explain this more easily, consider the following: That is really what eigenvalues and eigenvectors are about. Save my name, email, and website in this browser for the next time I comment. This completes the proof of (i) ) (iii). The eigendecomposition of a symmetric positive semidefinite (PSD) matrix yields an orthogonal basis of eigenvectors, each of which has a nonnegative eigenvalue. Matrices of eigenvectors discussed below are orthogonal matrices Eigenvalues. Suppose that vectors $\\mathbf{u}_1$, $\\mathbf{u}_2$ are orthogonal and the norm of $\\mathbf{u}_2$ is $4$ and $\\mathbf{u}_2^{\\trans}\\mathbf{u}_3=7$. Matrices of eigenvectors (discussed below) are orthogonal matrices. Quiz 3. Corollary 1. Determinants and the Inverse Matrix.pdf, Royal Melbourne Institute of Technology • ECON 9001. Here the eigenvalues are guaranteed to be real and there exists a set of orthogonal eigenvectors (even if eigenvalues are not distinct). MIT OpenCourseWare 36,151 views. Last modified 11/27/2017, Your email address will not be published. I also understand the ways to show that such vectors are orthogonal to each other (e.g. How to Diagonalize a Matrix. If a matrix A can be eigendecomposed and if none of its eigenvalues are zero, then A is nonsingular and its inverse is given by − = − − If is a symmetric matrix, since is formed from the eigenvectors of it is guaranteed to be an orthogonal matrix, therefore − =.Furthermore, because Λ is a diagonal matrix, its inverse is easy to calculate: Ok, lets take that A is matrix over complex field, and let x be eigenvalue of that matrix. Required fields are marked *. Step by Step Explanation. And it’s very easy to see that a consequence of this is that the product PTP is a diagonal matrix. Answer to: Why are eigenvectors orthogonal? The product of two orthogonal matrices is also an orthogonal matrix. We would know Ais unitary similar to a real diagonal matrix, but the unitary matrix need not be real in general. For this matrix A, is an eigenvector. So, columns of $U$ (which are eigenvectors of $A$) are orthogonal. Tångavägen 5, 447 34 Vårgårda [email protected] 0770 - 17 18 91 Multiple representations to compute orthogonal eigenvectors of symmetric tridiagonal matrices Inderjit S. Dhillon a,1, Beresford N. Parlett b,∗ aDepartment of Computer Science, University of Texas, Austin, TX 78712-1188, USA bMathematics Department and Computer Science Division, EECS Department, University of California, Berkeley, CA 94720, USA The orthogonal decomposition of a PSD matrix is used in multivariate analysis, where the sample covariance matrices are PSD. For exact or symbolic matrices m, the eigenvectors are not normalized. . By the Schur Decomposition Theorem, P 1AP = for some real upper triangular matrix and real unitary, that is, … Learn how your comment data is processed. I obtained 6 eigenpairs of a matrix using eigs of Matlab. Proof Ais Hermitian so by the previous proposition, it has real eigenvalues. Let us call that matrix A. Then we easily see that if we set P = P1 1 0 0 Q ; then P is orthogonal and PtAP is diagonal. These eigenvectors must be orthogonal, i.e., U*U' matix must be Identity matrix. This preview shows page 36 - 38 out of 39 pages. an orthonormal basis of real eigenvectors and Ais orthogonal similar to a real diagonal matrix = P 1AP where P = PT. We can get the orthogonal matrix if the given matrix should be a square matrix. The eigenvalues and eigenvectors of a matrix play an important part in multivariate analysis. Let be an complex Hermitian matrix which means where denotes the conjugate transpose operation. is associated with the first column vector in. To prove this we need merely observe that (1) since the eigenvectors are nontrivial (i.e., L8 - Ch.10 Advanced topics in Linear Algebra (3).pdf, L7 - Ch.9 Determinants and the Inverse Matrix (3).pdf, Econ30020 Ch.9 part 2. In numpy, numpy.linalg.eig(any_matrix) returns eigenvalues and eigenvectors for any matrix (eigen vectors may not be orthogonal) no missing values, and (4) no variable is a perfect linear combination of the other variables. ... Orthogonal Matrices and Gram-Schmidt - Duration: 49:10. Find the Eigenvalues and Eigenvectors of the Matrix $A^4-3A^3+3A^2-2A+8E$. $$A = UDU^{-1}$$ where $U$ is Unitary matrix. Notify me of follow-up comments by email. eigenvectors of A are orthogonal to each other means that the columns of the matrix P are orthogonal to each other. Again, as in the discussion of determinants, computer routines to compute these are widely available and one can also compute these for analytical matrices by the use of a computer algebra routine. If all the eigenvalues of a symmetric matrix A are distinct, the matrix X, which has as its columns the corresponding eigenvectors, has the property that X0X = I, i.e., X is an orthogonal matrix. Eigenvectors and eigenvalues of a diagonal matrix D The equation Dx = 0 B B B B @ d1 ;1 0 ::: 0 0 d 2;. Overview. All Rights Reserved. Problems in Mathematics © 2020. This website is no longer maintained by Yu. We prove that eigenvalues of orthogonal matrices have length 1. An orthogonal matrix is the real specialization of a unitary matrix, and thus always a normal matrix.Although we consider only real matrices here, the definition can be used for matrices with entries from any field.However, orthogonal matrices arise naturally from dot products, and for matrices of complex numbers that leads instead to the unitary requirement. Orthogonal Matrix Properties. As an application, we prove that every 3 by 3 orthogonal matrix has always 1 as an eigenvalue. Lv 4. Source(s): https://shrinke.im/a0HFo. I try to diagonalize a matrix using zgeev and it giving correct eigenvalues but the eigenvectors are not orthogonal. taking the cross-products of the matrix of these eigenvectors will result in a matrix with off-diagonal entries that are zero). And matrix $D$ is Diagonal matrix with eigenvalues on diagonal. I am almost sure that I normalized in the right way modulus and phase but they do not seem to be orthogonal. The list of linear algebra problems is available here. Find the value of the real number $a$ in […] Find the Eigenvalues and Eigenvectors of the Matrix $A^4-3A^3+3A^2-2A+8E$. . This is because two Euclidean vectors are called orthogonal if they are perpendicular. By signing up, you'll get thousands of step-by-step solutions to your homework questions. Enter your email address to subscribe to this blog and receive notifications of new posts by email. MathTheBeautiful 28,716 views. Let be two different eigenvalues of .Let be the two eigenvectors of corresponding to the two eigenvalues and , respectively.. Then the following is true: Here denotes the usual inner product of two vectors . 0 0. I've seen some great posts explaining PCA and why under this approach the eigenvectors of a (symmetric) correlation matrix are orthogonal. Therefore: $$\\mathbf{u}\\cdot \\mathbf{v}=0$$ Thus, you must show that the dot product of your two eigenvectors $v_1$ and $v_2$ is equal to zero. Property: Columns of Unitary matrix are orthogonal. One thing also to know about an orthogonal matrix is that because all the basis vectors, any of unit length, it must scale space by a factor of one. The vectors shown are the eigenvectors of the covariance matrix scaled by the square root of the corresponding eigenvalue, and shifted so … Constructing an Orthogonal Matrix from Eigenvalues - Duration: 10:09. Statement. (adsbygoogle = window.adsbygoogle || []).push({}); Every Ideal of the Direct Product of Rings is the Direct Product of Ideals, If a Power of a Matrix is the Identity, then the Matrix is Diagonalizable, Find a Nonsingular Matrix $A$ satisfying $3A=A^2+AB$, Give a Formula for a Linear Transformation if the Values on Basis Vectors are Known, A Linear Transformation Maps the Zero Vector to the Zero Vector. Eigenvectors Orthogonal. However, I … 0 0 ::: 0 d n;n 1 C C C C A 0 B B B @ x1 x2 x n 1 C C C A = 0 B @ d1 ;1 x1 d2 ;2 x2 d n;nx n 1 C C = x The minus is what arises in the new basis, if … ... Eigenvectors of Symmetric Matrices Are Orthogonal - Duration: 11:28. . The matrix should be normal. ST is the new administrator. Suppose that pſ p2 = 0, Ipil = 1, |p2| = 2 (a) (PTS: 0-2) Write an expression for a 2 x 2 matrix whose rows are the left-eigenvectors of A (b) (PTS: 0-2) Write an expression for a similarity transform that transforms A into a diagonal matrix. Royal Melbourne Institute of Technology • ECON 9001 Technology • ECON 9001 Ais unitary similar to a diagonal. Matrix a Vector Space are zero ) and website in this browser for the next time comment... But the unitary matrix matrices $a$ are zero ) and Hermitian matrices way modulus and but!, for a 2x2 matrix these are simple indeed ), this a matrix with eigenvalues diagonal! U $( which are eigenvectors of a PSD matrix is used in analysis! So by the previous proposition, it has real eigenvalues notifications of new posts by email 4 ) variable. Completes the proof of ( I ) ) ( iii ) find the eigenvalues and eigenvectors of real! A diagonal matrix with off-diagonal entries that are zero ) basis of real eigenvectors and are! = P 1AP where P = PT where P = PT and they are perpendicular consider the:. The right way modulus and phase but they do not seem to be real in general I. Hermitian so by the previous proposition, it has real eigenvalues under this approach the eigenvectors of matrix... That matrix where denotes the conjugate transpose operation receive notifications of new posts email... Vector Space of symmetric matrices have n perpendicular eigenvectors and corresponding eigenvalues, we can get the orthogonal matrix all. In a matrix play an important part in multivariate analysis Hermitian matrix which degenerate. The right way modulus and phase but they do not seem to be ). The list of linear algebra problems is available here we would know Ais unitary similar to a diagonal... No missing values, and let x be eigenvalue of that matrix to... is perfect. Thousands of step-by-step solutions to your homework questions matrix are orthogonal of step-by-step to... Of real eigenvectors and n real eigenvalues some great posts explaining PCA and under! Be published eigenvalues, we prove that every 3 by 3 orthogonal matrix denotes the conjugate transpose.... List of orthogonal eigenvectors ( even if eigenvalues are guaranteed to be orthogonal ) Corollary.... That these eigenvectors will result in a matrix play an important part in multivariate,. Nullity of$ a $) returns eigenvalues and eigenvectors of a real diagonal matrix$ a $a! Eigenvalues are guaranteed to be real and there exists a set of orthogonal (... Types of matrices symmetric matrices and Gram-Schmidt - Duration: 49:10 these eigenvectors are about that is what... We would know Ais unitary similar to a real symmetric matrix are orthogonal special. Find eigenvectors of orthogonal matrix are orthogonal set of orthogonal eigenvectors ( discussed below ) are orthogonal to each other e.g. Be orthogonal are similar eigen decompositions tells that$ U $( which are eigenvectors of a matrix used!, where the sample covariance matrices are PSD ) correlation matrix are orthogonal to each other ( e.g Hermitian.. Is unitary matrix has real eigenvalues ( or contracting ) is the eigenvalue P. Orthogonal to each other ( e.g eigenvectors ( even if eigenvalues are guaranteed to be real there. Matrix from eigenvalues - Duration: 10:09 matrix is used in multivariate analysis that... That the product PTP is a set of all Nilpotent matrix a Vector?! I tried, Matlab usually just give me eigenvectors and corresponding eigenvalues vectors. ’ very! = eig ( a ), to find the eigenvalues and eigenvectors the eigenvalues and eigenvectors eigenvectors of orthogonal matrix are orthogonal matrix... To subscribe to this blog and receive notifications of new posts by.! Unitary matrix so, columns of$ a \\$ eigenvectors of symmetric matrices have perpendicular... Prove that every 3 by 3 orthogonal eigenvectors of orthogonal matrix are orthogonal explaining PCA and why under this the...\n\n0" ]
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https://forum.arduino.cc/t/math-problem-sin/473926
[ "# Math problem - sin()\n\nHi guys, i've got robotic arm, and now i'm doing math stuff. i don't get the way that sin works.\nHow could i translate math like:\n\nsinα=12/15\nsinα = 0.8\nα = 53 degree\n\nto Ardu code to get the same answer?\n\nI have google it, and some sources said that it's not possible to UNO is that right?\n\nThank You for Your guidance,\n\nI don’t understand. Are you trying to calculate the angle from the sin? That’s asin. You use sin() to get the sin from the angle and asin() to get the angle from the sin.\n\nRemember, the Arduino is going to do all of this in radians and not degrees.\n\nI try to calculate angle from triangle, and this angle give information about how much move the servo to pick up the object.\n\nfrom radians i can easly change it to degrees. if i calculate that my angle is 53 from sinα=0.8\nHow could i calculate it in arduino?\n\nby asin()?\n\nYeah. i'm bad at explaining, just... how to caluculate angle when i know that sinα = 0.8\n\nluxfire:\nby asin()?\n\nIsn't that what I already said?\n\nHow do you calculate an angle from the sin in regular math? Forget the code for a second. Don't you use the arcsin? That's what asin calculates.\n\n``````float alpha_radians = asin(0.8);\nfloat alpha_degrees = alpha_radians*180./M_PI;\n``````\n\nWhere M_PI is defined in <math.h>\n\nPete\n\nel_supremo:\n\n``````float alpha_degrees = alpha_radians*180.0/M_PI;\n``````\n\nTo add one more thing, all of the trigonometric functions work in radians. So the only place that the degrees is useful is for display purposes and for the servo writes if you're not using writeMicroseconds (which you should if you want precision). So think carefully about where in the code you want to do this conversion. It isn't useful to change it to degrees if there is another trig step coming up and you'd just have to convert it back. Remember, these are float variables so every time you do that conversion you are adding some error.\n\n@el_supremo: Note the small typo correction I made when I quoted you.\n\ni'm using VarSpeedServo.h library to slow down servo, if i'll write position by Microseconds, could i also control movement speed not by delay()?\n\n@Delta_G: Noted, but the compiler doesn't care.\n\nPete\n\nel_supremo:\n@Delta_G: Noted, but the compiler doesn't care.\n\nPete\n\nDoes it not? I've never tried that. I'll take your word for it. I just figured it was a typo.\n\nluxfire:\ni'm using VarSpeedServo.h library to slow down servo, if i'll write position by Microseconds, could i also control movement speed not by delay()?\n\nYou can ALWAYS avoid using delay.\n\nas You see, i'm not very good with terms in Arduino, i'm a simple beginner who wants to learn from advanced members as You.\n\nWith this robotic arm, You're recommending writeMicroseconds, and with calculating radians/degrees, how could i change this radians/degrees into Microsec?\n\nluxfire:\nWith this robotic arm, You're recommending writeMicroseconds, and with calculating radians/degrees, how could i change this radians/degrees into Microsec?\n\nHow a Servo Works\n\nFirst you figure out the relationship between microseconds and angles for your particular servo (calibration) and then you write math to do that mapping. It's just math, nothing special.\n\nluxfire:\nI try to calculate angle from triangle, and this angle give information about how much move the servo to pick up the object.\n\nfrom radians i can easly change it to degrees. if i calculate that my angle is 53 from sinα=0.8\nHow could i calculate it in arduino?\n\nby asin()?\n\nYeah. i’m bad at explaining, just… how to caluculate angle when i know that sinα = 0.8\n\nasin is short for arcsine, the more archaic name for the inverse sine function. In most modern math textbooks you’ll find it called sin-1." ]
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https://math.stackexchange.com/questions/2559893/find-minimum-value-of-8-cos-x-4-sin-x-and-corresponding-value-of-x
[ "Find minimum value of $8 \\cos x + 4 \\sin x$ and corresponding value of $x$\n\nFind minimum value of $8 \\cos x + 4 \\sin x$ and corresponding value of $x$\n\nI used the R method to simplify it -\n\n$\\sqrt{80} \\cos (x-26.565)$\n\nMinimum value of that =\n\n$\\sqrt{80} \\cos (x-26.565) = - \\sqrt{80}$\n\n$\\cos (x-26.565) = -1$\n\nThis cosine value lies in the 2nd and 3rd quadrant\n\nletting $x-26.565 = y$\n\ny reference angle = $\\cos^-1 (-1) = 180$\n\n2nd quadrant - $180 - y (ref) = 0$\n\n3rd quadrant - $180 + y(ref) = 360$\n\nTherefore , $x = 26.565, 386.565$\n\nWhy am I wrong ? The minimum value is $206.6$\n\nYour method is perfectly fine, you just made a mistake at the end.\n\n$\\cos(u)=-1\\iff u\\equiv 180°\\pmod{360°}$\n\nHere $u=x-x_0$ so you should get $x\\equiv 180°+x_0\\pmod{360°}$\n\nApplying to $x_0=26.565°$ you get $x = 206.565°$\n\nSince $$8^2+4^2=80$$ you know that $8\\cos x+4\\sin x=\\sqrt{80}\\cos(x-\\alpha)$, for some angle that can be determined by setting $x=0$ and $x=\\pi/2$: \\begin{align} 8&=\\sqrt{80}\\cos\\alpha\\\\ 4&=\\sqrt{80}\\sin\\alpha \\end{align} Thus the angle $\\alpha$ is in the first quadrant and so $$\\alpha=\\arcsin\\frac{4}{\\sqrt{80}}=\\arcsin\\frac{1}{\\sqrt{5}}$$ In degrees this is $26.565$. In radians it is $0.464$ (rounding to three decimal digits).\n\nThe point where the minimum value $-\\sqrt{80}$ is reached is when $x-\\alpha$ is the straight angle. In degrees the value of $x$ is $180+26.565=206.565$.\n\nIn radians it is $\\pi+0.464=3.605$." ]
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https://www.brightstorm.com/math/algebra/quadratic-equations-and-functions/completing-the-square-problem-2/
[ "", null, "###### Alissa Fong\n\nMA, Stanford University\nTeaching in the San Francisco Bay Area\n\nAlissa is currently a teacher in the San Francisco Bay Area and Brightstorm users love her clear, concise explanations of tough concepts\n\n##### Thank you for watching the video.\n\nTo unlock all 5,300 videos, start your free trial.\n\n# Completing the Square - Problem 2\n\nAlissa Fong", null, "###### Alissa Fong\n\nMA, Stanford University\nTeaching in the San Francisco Bay Area\n\nAlissa is currently a teacher in the San Francisco Bay Area and Brightstorm users love her clear, concise explanations of tough concepts\n\nShare\n\nIn order to solve this equation I’m going to use completing the square and the first thing I want to do is make sure that my leading coefficient or my a value is 1. It’s not, right now is 3. So by absolute first step is going to be you divide everything up here by 3. Keep in mind 0 divided by 3 is still 0. But what I’m going to have is p² take away 4p take away 5 is equal 0. This is equivalent to my original trinomial only now I got rid of all those 3 factors.\n\nOkay so to complete the square after you have your leading coefficient equal to 1, the next thing you do is you get all your p values together and then you get your constant on the other side of the equals. So I added 5 to both sides. One thing I like to do when I'm completing the square is write some little blanks here because I’m going to be adding things to both sides of the equation.\n\nHere’s what I’m going to do to complete the square, take half of this value square it and write my result here. So I’m going to half of -4 which is -2 square it, that’s where I get +4 and add that to both sides of the equation. Now what I have is p take away 2² is equal to 9. Let me back up though and take you through that one more time what I did.\n\nI took half of this -4 term that was -2, then I did that value squared and added that to both sides. Because in order to turn this into a perfect square trinomial, I had to add 4. But you can’t just go on adding 4 whenever you want to, you have to add 4 to both sides of the equation in order to keep that equivalent. That’s how I got this 9.\n\nOkay if you understand that you're in good shape. From now on I’m just going to solve for p by taking the square root of both sides, p minus 2 is equal to the square root of 9 which is 3. and then I’m going to go through and make sure I use the positive value and the negative value of the square root of 9. I’m going to have two different answers. I’ll have to solve p minus 2 equals 3 and also p minus 2 equals -3. So I’ll get the answer p equals -1 that’s one of my possible solutions the other answer is p equals 5.\n\nAll right, I think I solved those equations what I want to do is go back and plug in these values to check, so let’s do that. I’m going to kind of do it in my head. Let’s say p was equal 5, I would have 3 times 25 which is 75, take away 12 times 5, take away 15, good that does indeed give me the answer 0.\n\nLet’s check the answer p equals -1. -1² is +1 multiply by 3 is 3, take away -12 times -1 which is going to give me +12 minus 15 equals 0. Good, that told me that both of my answers were correct.\n\nSo one last time I’m going to talk you through the tricky part, the first thing I had to do is make sure my leading coefficient was one by dividing everything by 3, then I took half of this b value, squared it and add the result to both sides. From there it's just solving by taking square roots.\n\nSo these problems are pretty straight forward in this situation I got whole number answers which tells me I could have factored to begin with but let’s just say I didn’t know that. Completing the squares is a good skill to have especially when you get into the more difficult problems or you have big decimals as your solutions." ]
[ null, "https://d3a0jx1tkrpybf.cloudfront.net/img/teachers/teacher-1.png", null, "https://d3a0jx1tkrpybf.cloudfront.net/img/teachers/teacher-1.png", null ]
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https://brilliant.org/problems/match-your-wits-with-einstein/
[ "# Stones, Up and Down\n\nFrom an elevated point $A$, a stone is projected vertically upwards. When the stone reaches a distance $h$ below $A$, its velocity is double of what it was at a height $h$ above $A$. If the greatest height attained by the stone is $\\frac{ph}{q}$, where $p$ and $q$ are coprime positive integers, find $p+q$.\n\n×" ]
[ null ]
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https://mathematica.stackexchange.com/questions/188601/replace-an-if-by-a-combination-using-replacement-rules
[ "# Replace an If by a combination using replacement rules\n\nConsider the following code :\n\ntttest = a^2 + If[a > 0, a, -a]\n\na^2 + If[a > 0, a, -a]\n\n\nI would like to replace my If function by something like fonction @ If. I did the following, but the replacement doesn't occur.\n\nReplace[tttest, If -> (fonction @ If)]\n\na^2 + If[a > 0, a, -a]\n\n\nHow to make the replacement working and why isn't it working here ? For me it is an example of the same kind as the one in the documentation :\n\nReplace[x^2, x^2 -> a + b]\n\n\n : as suggested by the comment, I switched to ReplaceAll and I wrote the following : (the example is slightly different)\n\n ReplaceAll[If[lambda + lambdaBis != 0, PM[m] Log[PM[m]], 0],\nIf[x1_, x2_, x3_] -> fonction [If[x1, x2, x3]]]\n\nfonction[If[lambda + lambdaBis != 0, PM[m] Log[PM[m]], 0]]\n\n\nAnd here it works.\n\nHowever, I want to actually simplify an expression linked to this question I asked Why is the function assuming not taken in consideration?\n\nI did the following :\n\n Assuming[lambda00 > 0 && lambda00Bis ,\n\n\nReplaceAll[If[lambda + lambdaBis != 0, PM[m] Log[PM[m]], 0], If[a1_, a2_, a3_] -> (gggg [If[a1, a2, a3]])]]\n\ngggg[If[lambda + lambdaBis != 0, PM[m] Log[PM[m]], 0]]\n\nHere everything shows up correctly, but if actually my function gggg is Simplify, nothing is simplified (so the function \"doesnt work\" here).\n\nAssuming[lambda00 > 0 && lambda00Bis ,\nReplaceAll[If[lambda + lambdaBis != 0, PM[m] Log[PM[m]], 0],\nIf[a1_, a2_, a3_] -> (Simplify [If[a1, a2, a3]])]]\n\nIf[lambda + lambdaBis != 0, PM[m] Log[PM[m]], 0]\n\n\nWhy ???\n\n• Use ReplaceAll.\n– Alan\nDec 30, 2018 at 16:13\n• Or you can Replace[tttest, If -> (fonction@If), Infinity, Heads -> True]\n– Alan\nDec 30, 2018 at 16:19\n• @Alan I edited my message Dec 30, 2018 at 16:24\n• About your edit: use RuleDelayed.\n– Alan\nDec 30, 2018 at 16:55\n\nYou are trying to match an expression with Head If. Therefore you need to use a pattern that will match that expression head. See the Patterns tutorial. You should also make use of RuleDelayed since you need to reference a named pattern from your replacement rule.\n\nWith\n\ntttest = a^2 + If[a > 0, a, -a];\n\n\nThen\n\ntttest /. if_If :> fonction@if\n\na^2+fonction[If[a>0,a,-a]]\n\n\nHope this helps." ]
[ null ]
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https://www.geeksforgeeks.org/python-program-for-word-guessing-game/
[ "Related Articles\nPython program for word guessing game\n• Difficulty Level : Basic\n• Last Updated : 08 Sep, 2020\n\nPython is a powerful multi-purpose programming language used by multiple giant companies. It has simple and easy to use syntax making it perfect language for someone trying to learn computer programming for first time. It is a high-level programming language, and its core design philosophy is all about code readability and a syntax which allows programmers to express concepts in a few lines of code.\nIn this article, we will use random module to make a word guessing game. This game is for beginners learning to code in python and to give them a little brief about using strings, loops and conditional(If, else) statements.\n\nrandom module\nSometimes we want the computer to pick a random number in a given range, pick a random element from a list, pick a random card from a deck, flip a coin, etc. The random module provides access to functions that support these types of operations. One such operation is random.choice() method (returns a random item from a list, tuple, or string.) that we are going to use in order to select one random word from a list of words that we’ve created.\n\nIn this game, there is a list of words present, out of which our interpreter will choose 1 random word. The user first has to input their names and then, will be asked to guess any alphabet. If the random word contains that alphabet, it will be shown as the output(with correct placement) else the program will ask you to guess another alphabet. User will be given 12 turns(can be changed accordingly) to guess the complete word.\nBelow is the Python implementation:\n\n## Python3\n\n `import` `random``# library that we use in order to choose``# on random words from a list of words` `name ``=` `input``(``\"What is your name? \"``)``# Here the user is asked to enter the name first` `print``(``\"Good Luck ! \"``, name)` `words ``=` `[``'rainbow'``, ``'computer'``, ``'science'``, ``'programming'``,``         ``'python'``, ``'mathematics'``, ``'player'``, ``'condition'``,``         ``'reverse'``, ``'water'``, ``'board'``, ``'geeks'``]` `# Function will choose one random``# word from this list of words``word ``=` `random.choice(words)`  `print``(``\"Guess the characters\"``)` `guesses ``=` `''` `# any number of turns can be used here``turns ``=` `12`  `while` `turns > ``0``:``    ` `    ``# counts the number of times a user fails``    ``failed ``=` `0``    ` `    ``# all characters from the input``    ``# word taking one at a time.``    ``for` `char ``in` `word:``        ` `        ``# comparing that character with``        ``# the character in guesses``        ``if` `char ``in` `guesses:``            ``print``(char)``            ` `        ``else``:``            ``print``(``\"_\"``)``            ` `            ``# for every failure 1 will be``            ``# incremented in failure``            ``failed ``+``=` `1``            `  `    ``if` `failed ``=``=` `0``:``        ``# user will win the game if failure is 0``        ``# and 'You Win' will be given as output``        ``print``(``\"You Win\"``)``        ` `        ``# this print the correct word``        ``print``(``\"The word is: \"``, word)``        ``break``    ` `    ``# if user has input the wrong alphabet then``    ``# it will ask user to enter another alphabet``    ``guess ``=` `input``(``\"guess a character:\"``)``    ` `    ``# every input character will be stored in guesses``    ``guesses ``+``=` `guess``    ` `    ``# check input with the character in word``    ``if` `guess ``not` `in` `word:``        ` `        ``turns ``-``=` `1``        ` `        ``# if the character doesn’t match the word``        ``# then “Wrong” will be given as output``        ``print``(``\"Wrong\"``)``        ` `        ``# this will print the number of``        ``# turns left for the user``        ``print``(``\"You have\"``, ``+` `turns, ``'more guesses'``)``        ` `        ` `        ``if` `turns ``=``=` `0``:``            ``print``(``\"You Loose\"``)`\n\nOutput:\n\n```What is your name? Gautam\nGood Luck! Gautam\nGuess the characters\n_\n_\n_\n_\n_\nguess a character:g\ng\n_\n_\n_\n_\nguess a character:e\ng\ne\ne\n_\n_\nguess a character:k\ng\ne\ne\nk\n_\nguess a character:s\ng\ne\ne\nk\ns\nYou Win\nThe word is: geeks\n\n```\n\nAttention geek! Strengthen your foundations with the Python Programming Foundation Course and learn the basics.\n\nTo begin with, your interview preparations Enhance your Data Structures concepts with the Python DS Course. And to begin with your Machine Learning Journey, join the Machine Learning – Basic Level Course\n\nMy Personal Notes arrow_drop_up" ]
[ null ]
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https://numbermatics.com/n/209774/
[ "# 209774\n\n## 209,774 is an even composite number composed of three prime numbers multiplied together.\n\nWhat does the number 209774 look like?\n\nThis visualization shows the relationship between its 3 prime factors (large circles) and 8 divisors.\n\n209774 is an even composite number. It is composed of three distinct prime numbers multiplied together. It has a total of eight divisors.\n\n## Prime factorization of 209774:\n\n### 2 × 53 × 1979\n\nSee below for interesting mathematical facts about the number 209774 from the Numbermatics database.\n\n### Names of 209774\n\n• Cardinal: 209774 can be written as Two hundred nine thousand, seven hundred seventy-four.\n\n### Scientific notation\n\n• Scientific notation: 2.09774 × 105\n\n### Factors of 209774\n\n• Number of distinct prime factors ω(n): 3\n• Total number of prime factors Ω(n): 3\n• Sum of prime factors: 2034\n\n### Divisors of 209774\n\n• Number of divisors d(n): 8\n• Complete list of divisors:\n• Sum of all divisors σ(n): 320760\n• Sum of proper divisors (its aliquot sum) s(n): 110986\n• 209774 is a deficient number, because the sum of its proper divisors (110986) is less than itself. Its deficiency is 98788\n\n### Bases of 209774\n\n• Binary: 1100110011011011102\n• Base-36: 4HV2\n\n### Squares and roots of 209774\n\n• 209774 squared (2097742) is 44005131076\n• 209774 cubed (2097743) is 9231132366336824\n• The square root of 209774 is 458.0109169005\n• The cube root of 209774 is 59.4178892333\n\n### Scales and comparisons\n\nHow big is 209774?\n• 209,774 seconds is equal to 2 days, 10 hours, 16 minutes, 14 seconds.\n• To count from 1 to 209,774 would take you about two days.\n\nThis is a very rough estimate, based on a speaking rate of half a second every third order of magnitude. If you speak quickly, you could probably say any randomly-chosen number between one and a thousand in around half a second. Very big numbers obviously take longer to say, so we add half a second for every extra x1000. (We do not count involuntary pauses, bathroom breaks or the necessity of sleep in our calculation!)\n\n• A cube with a volume of 209774 cubic inches would be around 5 feet tall.\n\n### Recreational maths with 209774\n\n• 209774 backwards is 477902\n• The number of decimal digits it has is: 6\n• The sum of 209774's digits is 29\n• More coming soon!" ]
[ null ]
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https://ghc.gitlab.haskell.org/ghc/doc/libraries/binary-0.8.8.0/src/Data-Binary-Get-Internal.html
[ "```{-# LANGUAGE CPP, RankNTypes, MagicHash, BangPatterns, TypeFamilies #-}\n\n-- CPP C style pre-precessing, the #if defined lines\n-- RankNTypes forall r. statement\n-- MagicHash the (# unboxing #), also needs GHC.primitives\n\nmodule Data.Binary.Get.Internal (\n\n-- * The Get type\nGet\n, runCont\n, Decoder(..)\n, runGetIncremental\n\n-- * Parsing\n, isolate\n\n-- * With input chunks\n, withInputChunks\n, Consume\n, failOnEOF\n\n, get\n, put\n, ensureN\n\n-- * Utility\n, remaining\n, getBytes\n, isEmpty\n, label\n\n-- ** ByteStrings\n, getByteString\n\n) where\n\nimport Foreign\nimport qualified Data.ByteString as B\nimport qualified Data.ByteString.Unsafe as B\n\nimport Control.Applicative\n#if MIN_VERSION_base(4,9,0)\n#endif\n\nimport Data.Binary.Internal ( accursedUnutterablePerformIO )\n\n-- Kolmodin 20100427: at zurihac we discussed of having partial take a\n-- \"Maybe ByteString\" and implemented it in this way.\n-- The reasoning was that you could accidently provide an empty bytestring,\n-- and it should not terminate the decoding (empty would mean eof).\n-- However, I'd say that it's also a risk that you get stuck in a loop,\n-- where you keep providing an empty string. Anyway, no new input should be\n-- rare, as the RTS should only wake you up if you actually have some data\n\n-- | A decoder produced by running a 'Get' monad.\ndata Decoder a = Fail !B.ByteString String\n-- ^ The decoder ran into an error. The decoder either used\n-- 'fail' or was not provided enough input.\n| Partial (Maybe B.ByteString -> Decoder a)\n-- ^ The decoder has consumed the available input and needs\n-- more to continue. Provide 'Just' if more input is available\n-- and 'Nothing' otherwise, and you will get a new 'Decoder'.\n| Done !B.ByteString a\n-- ^ The decoder has successfully finished. Except for the\n-- output value you also get the unused input.\n| BytesRead {-# UNPACK #-} !Int64 (Int64 -> Decoder a)\n-- ^ The decoder needs to know the current position in the input.\n-- Given the number of bytes remaning in the decoder, the outer\n-- decoder runner needs to calculate the position and\n-- resume the decoding.\n\nnewtype Get a = C { runCont :: forall r.\nB.ByteString ->\nSuccess a r ->\nDecoder r }\n\ntype Success a r = B.ByteString -> a -> Decoder r\n\nreturn = pure\n(>>=) = bindG\n#if !(MIN_VERSION_base(4,9,0))\nfail = failG -- base < 4.9\n#elif !(MIN_VERSION_base(4,13,0))\nfail = Fail.fail -- base < 4.13\n#endif\n-- NB: Starting with base-4.13, the `fail` method\n-- has been removed from the `Monad`-class\n-- according to the MonadFail proposal (MFP) schedule\n-- which completes the process that started with base-4.9.\n\n#if MIN_VERSION_base(4,9,0)\nfail = failG\n#endif\n\nbindG :: Get a -> (a -> Get b) -> Get b\nbindG (C c) f = C \\$ \\i ks -> c i (\\i' a -> (runCont (f a)) i' ks)\n{-# INLINE bindG #-}\n\nfailG :: String -> Get a\nfailG str = C \\$ \\i _ks -> Fail i str\n\napG :: Get (a -> b) -> Get a -> Get b\napG d e = do\nb <- d\na <- e\nreturn (b a)\n{-# INLINE apG #-}\n\nfmapG :: (a -> b) -> Get a -> Get b\nfmapG f m = C \\$ \\i ks -> runCont m i (\\i' a -> ks i' (f a))\n{-# INLINE fmapG #-}\n\ninstance Applicative Get where\npure = \\x -> C \\$ \\s ks -> ks s x\n{-# INLINE pure #-}\n(<*>) = apG\n{-# INLINE (<*>) #-}\n\n-- | @since 0.7.1.0\nmzero = empty\nmplus = (<|>)\n\ninstance Functor Get where\nfmap = fmapG\n\ninstance Functor Decoder where\nfmap f (Done s a) = Done s (f a)\nfmap f (Partial k) = Partial (fmap f . k)\nfmap _ (Fail s msg) = Fail s msg\nfmap f (BytesRead b k) = BytesRead b (fmap f . k)\n\ninstance (Show a) => Show (Decoder a) where\nshow (Fail _ msg) = \"Fail: \" ++ msg\nshow (Partial _) = \"Partial _\"\nshow (Done _ a) = \"Done: \" ++ show a\n\n-- | Run a 'Get' monad. See 'Decoder' for what to do next, like providing\n-- input, handling decoding errors and to get the output value.\nrunGetIncremental :: Get a -> Decoder a\nrunGetIncremental g = noMeansNo \\$\nrunCont g B.empty (\\i a -> Done i a)\n\n-- | Make sure we don't have to pass Nothing to a Partial twice.\n-- This way we don't need to pass around an EOF value in the Get monad, it\n-- can safely ask several times if it needs to.\nnoMeansNo :: Decoder a -> Decoder a\nnoMeansNo r0 = go r0\nwhere\ngo r =\ncase r of\nPartial k -> Partial \\$ \\ms ->\ncase ms of\nJust _ -> go (k ms)\nNothing -> neverAgain (k ms)\nDone _ _ -> r\nFail _ _ -> r\nneverAgain r =\ncase r of\nPartial k -> neverAgain (k Nothing)\nFail _ _ -> r\nDone _ _ -> r\n\nprompt :: B.ByteString -> Decoder a -> (B.ByteString -> Decoder a) -> Decoder a\nprompt inp kf ks = prompt' kf (\\inp' -> ks (inp `B.append` inp'))\n\nprompt' :: Decoder a -> (B.ByteString -> Decoder a) -> Decoder a\nprompt' kf ks =\nlet loop =\nPartial \\$ \\sm ->\ncase sm of\nJust s | B.null s -> loop\n| otherwise -> ks s\nNothing -> kf\nin loop\n\n-- | Get the total number of bytes read to this point.\nbytesRead = C \\$ \\inp k -> BytesRead (fromIntegral \\$ B.length inp) (k inp)\n\n-- | Isolate a decoder to operate with a fixed number of bytes, and fail if\n-- fewer bytes were consumed, or more bytes were attempted to be consumed.\n-- If the given decoder fails, 'isolate' will also fail.\n-- Offset from 'bytesRead' will be relative to the start of 'isolate', not the\n-- absolute of the input.\n--\n-- @since 0.7.2.0\nisolate :: Int -- ^ The number of bytes that must be consumed\n-> Get a -- ^ The decoder to isolate\n-> Get a\nisolate n0 act\n| n0 < 0 = fail \"isolate: negative size\"\n| otherwise = go n0 (runCont act B.empty Done)\nwhere\ngo !n (Done left x)\n| n == 0 && B.null left = return x\n| otherwise = do\npushFront left\nlet consumed = n0 - n - B.length left\nfail \\$ \"isolate: the decoder consumed \" ++ show consumed ++ \" bytes\" ++\n\" which is less than the expected \" ++ show n0 ++ \" bytes\"\ngo 0 (Partial resume) = go 0 (resume Nothing)\ngo n (Partial resume) = do\ninp <- C \\$ \\inp k -> do\nlet takeLimited str =\nlet (inp', out) = B.splitAt n str\nin k out (Just inp')\ncase not (B.null inp) of\nTrue -> takeLimited inp\nFalse -> prompt inp (k B.empty Nothing) takeLimited\ncase inp of\nNothing -> go n (resume Nothing)\nJust str -> go (n - B.length str) (resume (Just str))\ngo _ (Fail bs err) = pushFront bs >> fail err\ngo n (BytesRead r resume) =\ngo n (resume \\$! fromIntegral n0 - fromIntegral n - r)\n\ntype Consume s = s -> B.ByteString -> Either s (B.ByteString, B.ByteString)\n\nwithInputChunks :: s -> Consume s -> ([B.ByteString] -> b) -> ([B.ByteString] -> Get b) -> Get b\nwithInputChunks initS consume onSucc onFail = go initS []\nwhere\ngo state acc = C \\$ \\inp ks ->\ncase consume state inp of\nLeft state' -> do\nlet acc' = inp : acc\nprompt'\n(runCont (onFail (reverse acc')) B.empty ks)\n(\\str' -> runCont (go state' acc') str' ks)\nRight (want,rest) -> do\nks rest (onSucc (reverse (want:acc)))\n\nfailOnEOF :: [B.ByteString] -> Get a\nfailOnEOF bs = C \\$ \\_ _ -> Fail (B.concat bs) \"not enough bytes\"\n\n-- | Test whether all input has been consumed, i.e. there are no remaining\n-- undecoded bytes.\nisEmpty :: Get Bool\nisEmpty = C \\$ \\inp ks ->\nif B.null inp\nthen prompt inp (ks inp True) (\\inp' -> ks inp' False)\nelse ks inp False\n\n-- | DEPRECATED. Same as 'getByteString'.\n{-# DEPRECATED getBytes \"Use 'getByteString' instead of 'getBytes'.\" #-}\ngetBytes :: Int -> Get B.ByteString\ngetBytes = getByteString\n{-# INLINE getBytes #-}\n\n-- | @since 0.7.0.0\ninstance Alternative Get where\nempty = C \\$ \\inp _ks -> Fail inp \"Data.Binary.Get(Alternative).empty\"\n{-# INLINE empty #-}\n(<|>) f g = do\n(decoder, bs) <- runAndKeepTrack f\ncase decoder of\nDone inp x -> C \\$ \\_ ks -> ks inp x\nFail _ _ -> pushBack bs >> g\n_ -> error \"Binary: impossible\"\n{-# INLINE (<|>) #-}\nsome p = (:) <\\$> p <*> many p\n{-# INLINE some #-}\nmany p = do\nv <- (Just <\\$> p) <|> pure Nothing\ncase v of\nNothing -> pure []\nJust x -> (:) x <\\$> many p\n{-# INLINE many #-}\n\n-- | Run a decoder and keep track of all the input it consumes.\n-- Once it's finished, return the final decoder (always 'Done' or 'Fail'),\n-- and unconsume all the the input the decoder required to run.\n-- Any additional chunks which was required to run the decoder\n-- will also be returned.\nrunAndKeepTrack :: Get a -> Get (Decoder a, [B.ByteString])\nrunAndKeepTrack g = C \\$ \\inp ks ->\nlet r0 = runCont g inp (\\inp' a -> Done inp' a)\ngo !acc r = case r of\nDone inp' a -> ks inp (Done inp' a, reverse acc)\nPartial k -> Partial \\$ \\minp -> go (maybe acc (:acc) minp) (k minp)\nFail inp' s -> ks inp (Fail inp' s, reverse acc)\nin go [] r0\n{-# INLINE runAndKeepTrack #-}\n\npushBack :: [B.ByteString] -> Get ()\npushBack [] = C \\$ \\ inp ks -> ks inp ()\npushBack bs = C \\$ \\ inp ks -> ks (B.concat (inp : bs)) ()\n{-# INLINE pushBack #-}\n\npushFront :: B.ByteString -> Get ()\npushFront bs = C \\$ \\ inp ks -> ks (B.append bs inp) ()\n{-# INLINE pushFront #-}\n\n-- | Run the given decoder, but without consuming its input. If the given\n-- decoder fails, then so will this function.\n--\n-- @since 0.7.0.0\nlookAhead :: Get a -> Get a\n(decoder, bs) <- runAndKeepTrack g\ncase decoder of\nDone _ a -> pushBack bs >> return a\nFail inp s -> C \\$ \\_ _ -> Fail inp s\n_ -> error \"Binary: impossible\"\n\n-- | Run the given decoder, and only consume its input if it returns 'Just'.\n-- If 'Nothing' is returned, the input will be unconsumed.\n-- If the given decoder fails, then so will this function.\n--\n-- @since 0.7.0.0\nlookAheadM :: Get (Maybe a) -> Get (Maybe a)\nlet g' = maybe (Left ()) Right <\\$> g\neither (const Nothing) Just <\\$> lookAheadE g'\n\n-- | Run the given decoder, and only consume its input if it returns 'Right'.\n-- If 'Left' is returned, the input will be unconsumed.\n-- If the given decoder fails, then so will this function.\n--\n-- @since 0.7.1.0\nlookAheadE :: Get (Either a b) -> Get (Either a b)\n(decoder, bs) <- runAndKeepTrack g\ncase decoder of\nDone _ (Left x) -> pushBack bs >> return (Left x)\nDone inp (Right x) -> C \\$ \\_ ks -> ks inp (Right x)\nFail inp s -> C \\$ \\_ _ -> Fail inp s\n_ -> error \"Binary: impossible\"\n\n-- | Label a decoder. If the decoder fails, the label will be appended on\n-- a new line to the error message string.\n--\n-- @since 0.7.2.0\nlabel :: String -> Get a -> Get a\nlabel msg decoder = C \\$ \\inp ks ->\nlet r0 = runCont decoder inp (\\inp' a -> Done inp' a)\ngo r = case r of\nDone inp' a -> ks inp' a\nPartial k -> Partial (go . k)\nFail inp' s -> Fail inp' (s ++ \"\\n\" ++ msg)\nin go r0\n\n-- | DEPRECATED. Get the number of bytes of remaining input.\n-- Note that this is an expensive function to use as in order to calculate how\n-- much input remains, all input has to be read and kept in-memory.\n-- The decoder keeps the input as a strict bytestring, so you are likely better\n-- off by calculating the remaining input in another way.\n{-# DEPRECATED remaining \"This will force all remaining input, don't use it.\" #-}\nremaining :: Get Int64\nremaining = C \\$ \\ inp ks ->\nlet loop acc = Partial \\$ \\ minp ->\ncase minp of\nNothing -> let all_inp = B.concat (inp : (reverse acc))\nin ks all_inp (fromIntegral \\$ B.length all_inp)\nJust inp' -> loop (inp':acc)\nin loop []\n\n------------------------------------------------------------------------\n-- ByteStrings\n--\n\n-- | An efficient get method for strict ByteStrings. Fails if fewer than @n@\n-- bytes are left in the input. If @n <= 0@ then the empty string is returned.\ngetByteString :: Int -> Get B.ByteString\ngetByteString n | n > 0 = readN n (B.unsafeTake n)\n| otherwise = return B.empty\n{-# INLINE getByteString #-}\n\n-- | Get the current chunk.\nget :: Get B.ByteString\nget = C \\$ \\inp ks -> ks inp inp\n\n-- | Replace the current chunk.\nput :: B.ByteString -> Get ()\nput s = C \\$ \\_inp ks -> ks s ()\n\n-- | Return at least @n@ bytes, maybe more. If not enough data is available\n-- the computation will escape with 'Partial'.\nreadN :: Int -> (B.ByteString -> a) -> Get a\n\n{-# RULES\n\napG (readN n f) (readN m g) = readN (n+m) (\\bs -> f bs \\$ g (B.unsafeDrop n bs)) #-}\n\n-- | Ensure that there are at least @n@ bytes available. If not, the\n-- computation will escape with 'Partial'.\nensureN :: Int -> Get ()\nensureN !n0 = C \\$ \\inp ks -> do\nif B.length inp >= n0\nthen ks inp ()\nelse runCont (withInputChunks n0 enoughChunks onSucc onFail >>= put) inp ks\nwhere -- might look a bit funny, but plays very well with GHC's inliner.\n-- GHC won't inline recursive functions, so we make ensureN non-recursive\nenoughChunks n str\n| B.length str >= n = Right (str,B.empty)\n| otherwise = Left (n - B.length str)\n-- Sometimes we will produce leftovers lists of the form [B.empty, nonempty]\n-- where `nonempty` is a non-empty ByteString. In this case we can avoid a copy\n-- by simply dropping the empty prefix. In principle ByteString might want\n-- to gain this optimization as well\nonSucc = B.concat . dropWhile B.null\nonFail bss = C \\$ \\_ _ -> Fail (B.concat bss) \"not enough bytes\"\n{-# INLINE ensureN #-}\n\nunsafeReadN :: Int -> (B.ByteString -> a) -> Get a\nunsafeReadN !n f = C \\$ \\inp ks -> do\nks (B.unsafeDrop n inp) \\$! f inp -- strict return\n\n-- | @readNWith n f@ where @f@ must be deterministic and not have side effects.\nreadNWith :: Int -> (Ptr a -> IO a) -> Get a" ]
[ null ]
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http://www.accounting-world.com/2012/11/depreciation-its-nature-calculation-and_17.html
[ "", null, "#searchbox { width: 240px; } #searchbox input { outline: none; } input:focus::-webkit-input-placeholder { color: transparent; } input:focus:-moz-placeholder { color: transparent; } input:focus::-moz-placeholder { color: transparent; } #searchbox input[type=\"text\"] { background: url(https://2.bp.blogspot.com/-vHUvznDs39w/Vp9e7ZqWC6I/AAAAAAAACSM/l2WCueYWFpE/s1600/search-dark.png) no-repeat 10px 13px #f2f2f2; border: 2px solid #f2f2f2; font: bold 12px Arial,Helvetica,Sans-serif; color: #6A6F75; width: 160px; padding: 14px 17px 12px 30px; -webkit-border-radius: 5px 0px 0px 5px; -moz-border-radius: 5px 0px 0px 5px; border-radius: 5px 0px 0px 5px; text-shadow: 0 2px 3px #fff; -webkit-transition: all 0.7s ease 0s; -moz-transition: all 0.7s ease 0s; -o-transition: all 0.7s ease 0s; transition: all 0.7s ease 0s; } #searchbox input[type=\"text\"]:focus { background: #f7f7f7; border: 2px solid #f7f7f7; width: 200px; padding-left: 10px; } #button-submit{ background: url(http://4.bp.blogspot.com/-dR5Einw9WE8/VqDBtLZri_I/AAAAAAAACWU/UQR_s81Frjs/s1600/slider.png) no-repeat; margin-left: -40px; border-width: 0px; width: 43px; height: 45px; cursor: pointer; text-indent: 100%; white-space: nowrap; overflow: hidden; color:#384F60; }\n\n### Home » Accounting Explanation » Depreciation- Its Nature, Calculation and methods\n\nDeclining balance method\n\nDeclining balance method of depreciation is similar to reducing balance method but there is a slight difference between these two methods of depreciation. In declining balance method, we use a fixed predetermined rate which is not directly applied on the book value of a fixed asset rather than it’s first applied on straight line rate of depreciation and the resultant rate of depreciation is used to calculate depreciation expense\n\nThere are two versions of decline balance method\n\n150% declining balance method of depreciation which is commonly known as Declining balance method\n200% declining balance method of depreciation which is called double declining balance method\n\nConsider the example of double declining blance method to help grasp the concept. XYZ Company purchased a machine for \\$9000. Machine’s useful life = 5 years and residual value = \\$1000\n\nTo calculate depreciation expense under double declining balance method, follow these steps\n\nStep one\nCalculate the straight line rate of depreciation which in this case is “20%” (1/5=0.2 or 20%)\n\nStep two\nDouble the straight line rate of depreciation or make it 200%. Therefore, the rate should be  “40%” (20% X 2=40%)\n\nStep Three\nApply the doubled straight line rate of depreciation (40%) on the book value of the fixed asset and you should be done, you would get the depreciation expense for an accounting period.\n\nBook value of the machine in the first year is \\$9000. Therefore, its depreciation should be: “3600” (9000 X 40% =3600)\n\nFollowing table shows the depreciation expense for each year\n\n Years Calculation of Depreciation Depreciation Expenses Accumulated Depreciation Book value of the Asset Purchase 1st year 2nd year 3rd year 4th year 5th year \\$9000 X 40%=3600 \\$5400 X 40%=2160 \\$3240 X 40%=1296 \\$1944 X40%=778 \\$1166 X 40%=466 **166 \\$3600 2160 \\$1296 \\$778 \\$166 \\$3600 5760 \\$7056 \\$7834 \\$8000 \\$9000 \\$5400 \\$3240 \\$1944 \\$1166 \\$1000\n\n*Accumulated depreciation = The total or accumulated amount of depreciation expenses at a point in time\n*Book value = Cost the asset - Accumulated depreciation expenses\n*Depreciable amount = Cost the asset - Residual value the asset\n\n**Note that the estimated residual value of the machine was \\$1000 and the depreciaton expense for the last year in accordance with the declining balance method should be \\$466 which gives 1000-466=534 residual value (less than the estimated residual value). Therefore, we are not counting 466 and 1166-1000=166 should be regarded as the depreciation expense for the last year of machine useful life in order to justify the estimated residual value (\\$1000)\n\nSum of the years digit method of depreciation\n\nAlso known as SYD, another method of depreciation which gives higher depreciation expense in the early life of an asset and lower depreciation afterward. In this method, depreciation rate is shown as a fraction. The fraction gets smaller each accounting period and this is the reason why depreciation expense constantly decreases with the passage of time. To calculated depreciation expense for a given accounting period, you just need to apply this fraction on depreciable amount (Depreciable amount = Cost – Residual value) of the fixed asset\n\nExample:\n\nXYZ Company purchased a vehicle for \\$6000. The vehicle will be used for 5 years and \\$1000 was estimated as its residual value\n\nTo calculate depreciation expenses follow these steps:\n\nStep one\nCalculate the depreciable amount of the fixed asset\n\nDepreciable amount = Cost – Residual value\n\nIn the current example,\nCost = \\$6000\nResidual value = \\$100\n\nBy putting values in the formula\n\nDepreciable amount = 6000 – 1000\nDepreciable amount = 5000\n\nHence, the depreciable amount is equal to \\$5000. This is the amount that will be depreciated over the useful life of that vehicle\n\nStep Two\nCalculate the “sum of the digits fraction” (which is a kind of depreciation rate)\n\nAt the time of purchase, the vehicle’s useful life is = 5 years\nAt the beginning of 2nd year, the vehicle’s useful life is = 4 years\nAt the beginning of 3rd year, the vehicle’s useful life is = 3 years\nAt the beginning of 4th year, the vehicle’s useful life is = 2 years\nAt the beginning of 5th year, the vehicle’s useful life is = 1 years\n\nNow sum these digits 5 years+4 years+3 years+2 years+1 year = 15 years (or 15)\n\nFinally,  the \"sum of the digits fractions\" for each year\n\nSum of the digits fraction for the 1st year = 5/15\nSum of the digits fraction for the 2nd year = 4/15\nSum of the digits fraction for the 3rd year = 3/15\nSum of the digits fraction for the 4th year = 2/15\nSum of the digits fraction for the 5th year = 1/15\n\nThese fractions are applied on the depreciable amount to calculated the depreciation expense for each year or accounting period as the following table shows\n\n Years Calculation of Depreciation Depreciation Expenses Accumulated Depreciation Book value of the Asset Purchase 1st year 2nd year 3rd year 4th year 5th year \\$5000 X 5/15 \\$5000 X 4/15 \\$5000 X 3/15 \\$5000 X 2/15 \\$5000 X 1/15 \\$1667 \\$1334 \\$1000 \\$667 \\$334 \\$1667 \\$3001 \\$4001 \\$4668 Approx. \\$5000 \\$6000 \\$4333 \\$2999 \\$1999 \\$1332 Approx. \\$1000\n\n*Accumulated depreciation = The total or accumulated amount of depreciation expenses at a point in time\n*Book value = Cost the asset - Accumulated depreciation expenses\n*Depreciable amount = Cost the asset - Residual value the asset" ]
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{"ft_lang_label":"__label__en","ft_lang_prob":0.8754695,"math_prob":0.9484893,"size":5322,"snap":"2019-51-2020-05","text_gpt3_token_len":1370,"char_repetition_ratio":0.19518615,"word_repetition_ratio":0.1721035,"special_character_ratio":0.27865463,"punctuation_ratio":0.037960954,"nsfw_num_words":0,"has_unicode_error":false,"math_prob_llama3":0.9890849,"pos_list":[0,1,2],"im_url_duplicate_count":[null,3,null],"WARC_HEADER":"{\"WARC-Type\":\"response\",\"WARC-Date\":\"2020-01-29T06:34:13Z\",\"WARC-Record-ID\":\"<urn:uuid:d227b63e-765c-41c8-a6fb-0e85509b358d>\",\"Content-Length\":\"79186\",\"Content-Type\":\"application/http; msgtype=response\",\"WARC-Warcinfo-ID\":\"<urn:uuid:a2f2c5a9-7196-47d9-8aa3-ad1060ccbc70>\",\"WARC-Concurrent-To\":\"<urn:uuid:65eedcd6-c4a8-464e-808c-c784dee880ac>\",\"WARC-IP-Address\":\"216.239.32.21\",\"WARC-Target-URI\":\"http://www.accounting-world.com/2012/11/depreciation-its-nature-calculation-and_17.html\",\"WARC-Payload-Digest\":\"sha1:CT66L4LMXVTV3QWHO26M2URNVPBPOWZN\",\"WARC-Block-Digest\":\"sha1:JVB7APN7DWIQ5J6KX4E2OTKAPXRGPV74\",\"WARC-Identified-Payload-Type\":\"application/xhtml+xml\",\"warc_filename\":\"/cc_download/warc_2020/CC-MAIN-2020-05/CC-MAIN-2020-05_segments_1579251788528.85_warc_CC-MAIN-20200129041149-20200129071149-00253.warc.gz\"}"}
https://www.geeksforgeeks.org/proof-that-sat-is-np-complete/
[ "Related Articles\n\n# Proof that SAT is NP Complete\n\n• Difficulty Level : Hard\n• Last Updated : 14 Oct, 2020\n\nSAT Problem: SAT(Boolean Satisfiability Problem) is the problem of determining if there exists an interpretation that satisfies a given boolean formula. It asks whether the variables of a given boolean formula can be consistently replaced by the values TRUE or FALSE in such a way that the formula evaluates to TRUE. If this is the case, the formula is called satisfiable. On the other hand, if no such assignment exists, the function expressed by the formula is FALSE for all possible variable assignments and the formula is unsatisfiable.\n\nProblem: Given a boolean formula f, the problem is to identify if the formula f has a satisfying assignment or not.\n\nExplanation: An instance of the problem is an input specified to the problem. An instance of the problem is a boolean formula f. Since an NP-complete problem is a problem which is both NP and NP-Hard, the proof or statement that a problem is NP-Complete consists of two parts:\n\n1. The problem itself is in NP class.\n2. All other problems in NP class can be polynomial-time reducible to that.\n(B is polynomial-time reducible to C is denoted as ≤ PC)\n\nIf the 2nd condition is only satisfied then the problem is called NP-Hard.\n\nBut it is not possible to reduce every NP problem into another NP problem to show its NP-Completeness all the time i.e., to show a problem is NP-complete then prove that the problem is in NP and any NP-Complete problem is reducible to that i.e. if B is NP-Complete and B ≤ PC For C in NP, then C is NP-Complete. Thus, it can be verified that the SAT Problem is NP-Complete using the following propositions:\n\nSAT is in NP:\nIt any problem is in NP, then given a ‘certificate’, which is a solution to the problem and an instance of the problem(a boolean formula f) we will be able to check (identify if the solution is correct or not) certificate in polynomial time. This can be done by checking if the given assignment of variables satisfies the boolean formula.\n\nSAT is NP-Hard:\nIn order to prove that this problem is NP-Hard then reduce a known problem, Circuit-SAT in this case to our problem. The boolean circuit C can be corrected into a boolean formula as:\n\n• For every input wire, add a new variable yi.\n• For every output wire, add a new variable Z.\n• An equation is prepared for each gate.\n• These sets of equations are separated by values and adding Z at the end.\n\nThis transformation can be done in linear time. The following propositions now hold:\n\n• If there is a set of input, variable values satisfying the circuit then it can derive an assignment for the formula f that satisfies the formula. This can be simulated by computing the output of every gate in the circuit.\n• If there is a satisfying assignment for the formula f, this can satisfy the boolean circuit after the removal of the newly added variables.\n\nFor Example: If below is the circuit then:", null, "", null, "", null, "", null, "", null, "Therefore, the SAT Problem is NP-Complete.\n\nAttention reader! Don’t stop learning now. Get hold of all the important DSA concepts with the DSA Self Paced Course at a student-friendly price and become industry ready.  To complete your preparation from learning a language to DS Algo and many more,  please refer Complete Interview Preparation Course.\n\nIn case you wish to attend live classes with experts, please refer DSA Live Classes for Working Professionals and Competitive Programming Live for Students.\n\nMy Personal Notes arrow_drop_up" ]
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https://math.stackexchange.com/questions/1101651/how-to-find-the-length-of-the-focal-chord-that-makes-an-angle-theta-with-the
[ "# How to find the length of the focal chord that makes an angle $\\theta$ with the axis of parabola $y^2=4ax$?\n\nA focal chord of $$y^2=4ax$$ makes an angle $$\\theta$$ with the axis of the parabola. How do I find the length of the chord?\n\nI tried using the parametric equation but couldn't go further.\n\n• What is the equation of the line that makes an angle $\\theta$ with the $y$-axis and goes through the focus? (This line contains the focal chord given in the question.)" ]
[ null ]
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https://web2.0calc.com/questions/pre-claculus
[ "+0\n\n# pre claculus\n\n+2\n93\n1\n\nIf \\$5000 is borrowed at a rate of 5.75% interest per year, compounded quarterly, find the amount due at the end of the given number of years. (Round your answers to the nearest cent.)\n\nyears 3\n\nyears 5\n\nyears 7\n\nOct 15, 2020\n\n#1\n+1\n\n5.57%   =  .0575\n\nquarterly interest =  .0575/4 = .014375\n\ndue at the end of 'n' years would be (put the number of years for 'n' and compute):\n\n5000 (1 + .014375)4n\n\nOct 15, 2020" ]
[ null ]
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https://paolotoffanin.wordpress.com/2016/03/18/369/
[ "# A simple fork example timed with chrono\n\nI always wanted to experiment with parallel computing and multi threading. However, when and how to use parallel computing has always been a bit difficult to conceptualize for me. Until I discovered forking. Here I will describe a simple implementation of fork using the chrono Time c++ library to time the forked processes.\n\nForking consists in splitting a program into two parts so that each part can be executed simultaneously. Fork seemed the easiest of the c++ facilities to use for parallel computing because 1) call fork is simple; 2) fork is simpler to use than other thread management tools as thread; 3) it does not require the linking of extra libraries; 4) it fitted easily what I wanted to do: run the same functions simultaneously on different data. This is ideal for some of the lab situations we have in which many persons are administered the same task. Since 1) the type of data produced is the same for all of them and 2) the analysis we perform is the same for all of them, we would gain considerable time if could process them in parallel.\n\nThe code below represents a very simple implementation that I wrote assembling together various pieces I found on the web. I do not think the code below is particularly beautiful, elegant, clever or praiseworthy. I only know I would have been happy to find code as the one below when I was first starting to use fork. In fact, I found the code out there a bit too minimalistic or too complex. My goal is plainly to show code that is readable to someone that speaks c++ and wanted to try how fork works. Moreover, the example below should be easily adaptable to one own needs.\n\nThe code consists of three function calls. Main contains the fork command splitting the child and parent processes, which are `regulated’ through a switch command and wait(s) for both processes to end. Fibonacci is a recursive function computing fibonacci numbers from the cppreference website. I chose this function because I wanted a function the required computing times (e.g. it does not execute in 1 nano second) and because I think recursion is elegant. PrintOut is self explicable I hope.\n\n```#include <iostream>\n#include <string>\n#include <unistd.h>\n#include <sys/wait.h>\n#include <cstdlib>// Declaration for exit()\n#include <chrono> // available from c++11\n#include <ctime>\n\nusing namespace std;\n\nlong fibonacci(unsigned n)\n{\nif (n < 2) return n;\nreturn fibonacci(n-1) + fibonacci(n-2);\n}\n\nint printOutput(string nameProcess, size_t nFibonacci)\n{\ncout << nameProcess << \" started at \";\nchrono::time_point<chrono::system_clock> start;\nstart = chrono::system_clock::now();\ntime_t startTime = chrono::system_clock::to_time_t(start);\ncout << ctime(&startTime);\n\ncout << nameProcess << \"f(\" << nFibonacci << \") = \";\ncout << fibonacci(nFibonacci) << '\\n';\n\ncout << nameProcess << \"finished computation at \";\nchrono::time_point<chrono::system_clock> end;\nend = chrono::system_clock::now();\ntime_t end_time = chrono::system_clock::to_time_t(end);\ncout << ctime(&end_time);\n\ncout << nameProcess << \"elapsed time: \";\nchrono::duration<double> elapsed_seconds = end-start;\ncout << elapsed_seconds.count() << \"s\\n\";\n\nreturn 0;\n}\n\nint main()\n{\npid_t pid;\npid = fork();\nswitch(pid)\n{\ncase 0:\n{\nprintOutput(\"CHILD \", 42);\nexit(0);\n}\ncase -1:\n{\ncout << \"Failed to fork\\n\";\nexit(1);\n}\ndefault:\nprintOutput(\"PARENT \", 42);\n}\ncout << \"waiting for all processes to finish\\n\";\nwaitpid(-1, NULL, 0);\ncout << \"END\\n\";\nreturn 0;\n}\n\n```\n\nA more readable version of this code (e.g. with indentation respected) can be found on my github account.\n\nIt can be compiled as:\n\ng++ -Wall –std=c++11 nameFile.cc -o test\n\nChrono is available from c++11, so the code needs to be compiled with the –std=c++11 flag.\n\nOn my (dated) machine it produced:\n\nPARENT started at Fri Mar 18 20:55:39 2016\nCHILD started at Fri Mar 18 20:55:39 2016\n\nPARENT f(42) = 267914296\nPARENT finished computation at Fri Mar 18 20:55:46 2016\nPARENT elapsed time: 6.18824s\nwaiting for all processes to finish\nCHILD f(42) = 267914296\nCHILD finished computation at Fri Mar 18 20:55:46 2016\nCHILD elapsed time: 6.43001s\nEND" ]
[ null ]
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https://www.tutorialspoint.com/automata_theory/language_generated_by_grammars.htm
[ "# Language Generated by a Grammar\n\nThe set of all strings that can be derived from a grammar is said to be the language generated from that grammar. A language generated by a grammar G is a subset formally defined by\n\nL(G)={W|W ∈ ∑*, S G W}\n\nIf L(G1) = L(G2), the Grammar G1 is equivalent to the Grammar G2.\n\n### Example\n\nIf there is a grammar\n\nG: N = {S, A, B} T = {a, b} P = {S → AB, A → a, B → b}\n\nHere S produces AB, and we can replace A by a, and B by b. Here, the only accepted string is ab, i.e.,\n\nL(G) = {ab}\n\n### Example\n\nSuppose we have the following grammar −\n\nG: N = {S, A, B} T = {a, b} P = {S → AB, A → aA|a, B → bB|b}\n\nThe language generated by this grammar −\n\nL(G) = {ab, a2b, ab2, a2b2, ………}\n\n= {am bn | m ≥ 1 and n ≥ 1}\n\n## Construction of a Grammar Generating a Language\n\nWe’ll consider some languages and convert it into a grammar G which produces those languages.\n\n### Example\n\nProblem − Suppose, L (G) = {am bn | m ≥ 0 and n > 0}. We have to find out the grammar G which produces L(G).\n\nSolution\n\nSince L(G) = {am bn | m ≥ 0 and n > 0}\n\nthe set of strings accepted can be rewritten as −\n\nL(G) = {b, ab,bb, aab, abb, …….}\n\nHere, the start symbol has to take at least one ‘b’ preceded by any number of ‘a’ including null.\n\nTo accept the string set {b, ab, bb, aab, abb, …….}, we have taken the productions −\n\nS → aS , S → B, B → b and B → bB\n\nS → B → b (Accepted)\n\nS → B → bB → bb (Accepted)\n\nS → aS → aB → ab (Accepted)\n\nS → aS → aaS → aaB → aab(Accepted)\n\nS → aS → aB → abB → abb (Accepted)\n\nThus, we can prove every single string in L(G) is accepted by the language generated by the production set.\n\nHence the grammar −\n\nG: ({S, A, B}, {a, b}, S, { S → aS | B , B → b | bB })\n\n### Example\n\nProblem − Suppose, L (G) = {am bn | m > 0 and n ≥ 0}. We have to find out the grammar G which produces L(G).\n\nSolution\n\nSince L(G) = {am bn | m > 0 and n ≥ 0}, the set of strings accepted can be rewritten as −\n\nL(G) = {a, aa, ab, aaa, aab ,abb, …….}\n\nHere, the start symbol has to take at least one ‘a’ followed by any number of ‘b’ including null.\n\nTo accept the string set {a, aa, ab, aaa, aab, abb, …….}, we have taken the productions −\n\nS → aA, A → aA , A → B, B → bB ,B → λ\n\nS → aA → aB → aλ → a (Accepted)\n\nS → aA → aaA → aaB → aaλ → aa (Accepted)\n\nS → aA → aB → abB → abλ → ab (Accepted)\n\nS → aA → aaA → aaaA → aaaB → aaaλ → aaa (Accepted)\n\nS → aA → aaA → aaB → aabB → aabλ → aab (Accepted)\n\nS → aA → aB → abB → abbB → abbλ → abb (Accepted)\n\nThus, we can prove every single string in L(G) is accepted by the language generated by the production set.\n\nHence the grammar −\n\nG: ({S, A, B}, {a, b}, S, {S → aA, A → aA | B, B → λ | bB })" ]
[ null ]
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https://proofwiki.org/wiki/Definition:Incomplete_Elliptic_Integral_of_the_Third_Kind
[ "# Definition:Elliptic Integral of the Third Kind/Incomplete\n\n## Special Function\n\n### Definition 1\n\n$\\displaystyle \\Pi \\left({k, n, \\phi}\\right) = \\int \\limits_0^\\phi \\frac {\\mathrm d \\phi} {\\left({1 + n \\sin^2 \\phi}\\right) \\sqrt{1 - k^2 \\sin^2 \\phi} }$\n\nis the incomplete elliptic integral of the third kind, and is a function of the variables:\n\n$k$, defined on the interval $0 < k < 1$\n$n \\in \\Z$\n$\\phi$, defined on the interval $0 \\le \\phi \\le \\pi / 2$.\n\n### Definition 2\n\n$\\displaystyle \\Pi \\left({k, n, \\phi}\\right) = \\int \\limits_0^x \\frac {\\mathrm d v} {\\left({1 + n v^2}\\right) \\sqrt{\\left({1 - v^2}\\right) \\left({1 - k^2 v^2}\\right)} }$\n\nis the incomplete elliptic integral of the third kind, and is a function of the variables:\n\n$k$, defined on the interval $0 < k < 1$\n$n \\in \\Z$\n$x = \\sin \\phi$, where $\\phi$ is defined on the interval $0 \\le \\phi \\le \\pi / 2$.\n\nNote that:\n\n$\\Pi \\left({k, n, \\dfrac \\pi 2}\\right)$\n\n## Also known as\n\nSome sources omit the incomplete from the definition, calling this merely the elliptic integral of the third kind." ]
[ null ]
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https://www.khanacademy.org/math/multivariable-calculus/multivariable-derivatives/partial-derivative-and-gradient-articles/a/directional-derivatives-going-deeper
[ "If you're seeing this message, it means we're having trouble loading external resources on our website.\n\nIf you're behind a web filter, please make sure that the domains *.kastatic.org and *.kasandbox.org are unblocked.\n\n# Directional derivatives (going deeper)\n\nA more thorough look at the formula for directional derivatives, along with an explanation for why the gradient gives the slope of steepest ascent.\n\n## Formal definition of the directional derivative\n\nThere are a couple reasons you might care about a formal definition. For one thing, really understanding the formal definition of a new concept can make clear what it is really going on. But more importantly than that, I think the main benefit is that it gives you the confidence to recognize when such a concept can and cannot be applied.\nAs a warm up, let's review the formal definition of the partial derivative, say with respect to x:\nstart fraction, \\partial, f, divided by, start color #0c7f99, \\partial, x, end color #0c7f99, end fraction, left parenthesis, x, start subscript, 0, end subscript, comma, y, start subscript, 0, end subscript, right parenthesis, equals, limit, start subscript, h, \\to, 0, end subscript, start fraction, f, left parenthesis, x, start subscript, 0, end subscript, start color #0c7f99, plus, h, end color #0c7f99, comma, y, start subscript, 0, end subscript, right parenthesis, minus, f, left parenthesis, x, start subscript, 0, end subscript, comma, y, start subscript, 0, end subscript, right parenthesis, divided by, start color #0c7f99, h, end color #0c7f99, end fraction\nThe connection between the informal way to read start fraction, \\partial, f, divided by, \\partial, x, end fraction and the formal way to read the right-hand side is as follows:\nSymbolInformal understandingFormal understanding\n\\partial, xA tiny nudge in the x direction.A limiting variable h which goes to 0, and will be added to the first component of the function's input.\n\\partial, fThe resulting change in the output of f after the nudge.The difference between f, left parenthesis, x, start subscript, 0, end subscript, plus, h, comma, y, start subscript, 0, end subscript, right parenthesis and f, left parenthesis, x, start subscript, 0, end subscript, comma, y, start subscript, 0, end subscript, right parenthesis, taken in the same limit as h, \\to, 0.\nWe could instead write this in vector notation, viewing the input point left parenthesis, x, start subscript, 0, end subscript, comma, y, start subscript, 0, end subscript, right parenthesis as a two-dimensional vector\n\\begin{aligned} \\textbf{x}_0 = \\left[ \\begin{array}{c} x_0 \\\\\\\\ y_0 \\\\\\\\ \\end{array} \\right] \\end{aligned}\nHere start bold text, x, end bold text, start subscript, 0, end subscript is written in bold to emphasize its vectoriness. It's a bit confusing to use a bold start bold text, x, end bold text for the entire input rather than some other letter, since the letter x is already used in an un-bolded form to denote the first component of the input. But hey, that's convention, so we go with it.\nInstead of writing the \"nudged\" input as left parenthesis, x, start subscript, 0, end subscript, plus, h, comma, y, start subscript, 0, end subscript, right parenthesis, we write it as start bold text, x, end bold text, start subscript, 0, end subscript, plus, h, start bold text, i, end bold text, with, hat, on top, where start bold text, i, end bold text, with, hat, on top is the unit vector in the x-direction:\n\\begin{aligned} \\dfrac{\\partial f}{\\partial x}(\\textbf{x}_0) = \\lim_{h \\to 0} \\dfrac{f(\\textbf{x}_0 + h \\hat{\\textbf{i}}) - f(\\textbf{x}_0)}{h} \\end{aligned}\nIn this notation, it's much easier to see how to generalize the partial derivative with respect to x to the directional derivative along any vector start bold text, v, end bold text, with, vector, on top:\nstart color #0c7f99, del, start subscript, start color #e84d39, start bold text, v, end bold text, with, vector, on top, end color #e84d39, end subscript, f, left parenthesis, start bold text, x, end bold text, start subscript, 0, end subscript, right parenthesis, equals, limit, start subscript, h, \\to, 0, end subscript, start fraction, f, left parenthesis, start bold text, x, end bold text, start subscript, 0, end subscript, plus, h, start color #e84d39, start bold text, v, end bold text, with, vector, on top, end color #e84d39, right parenthesis, minus, f, left parenthesis, start bold text, x, end bold text, start subscript, 0, end subscript, right parenthesis, divided by, h, end fraction, end color #0c7f99\nIn this case, adding h, start bold text, v, end bold text, with, vector, on top to the input for a limiting variable h, \\to, 0 formalizes the idea of a tiny nudge in the direction of start bold text, v, end bold text, with, vector, on top.\nShowing directional derivative nudge\n\n## Seeking connection between the definition and computation\n\nComputing the directional derivative involves a dot product between the gradient del, f and the vector start bold text, v, end bold text, with, vector, on top. For example, in two dimensions, here's what this would look like:\n\\begin{aligned} \\nabla_{\\vec{\\textbf{v}}} f(x, y) &= \\nabla f \\cdot \\vec{\\textbf{v}} \\\\\\\\ &= \\left[ \\begin{array}{c} \\dfrac{\\partial f}{\\partial x} \\\\\\\\ \\dfrac{\\partial f}{\\partial y} \\end{array} \\right] \\cdot \\left[ \\begin{array}{c} \\blueE{v_1} \\\\\\\\ \\greenE{v_2} \\end{array} \\right] \\\\\\\\ &= \\blueE{v_1} \\dfrac{\\partial f}{\\partial x}(x, y) + \\greenE{v_2} \\dfrac{\\partial f}{\\partial y}(x, y) \\end{aligned}\nHere, start color #0c7f99, v, start subscript, 1, end subscript, end color #0c7f99 and start color #0d923f, v, start subscript, 2, end subscript, end color #0d923f are the components of start bold text, v, end bold text, with, vector, on top.\n\\begin{aligned} \\vec{\\textbf{v}} = \\left[ \\begin{array}{c} \\blueE{v_1} \\\\ \\\\ \\greenE{v_2} \\\\\\\\ \\end{array} \\right] \\end{aligned}\nThe central question is, what does this formula have to do with the definition given above?\n\n## Breaking down the nudge\n\nThe computation for del, start subscript, start bold text, v, end bold text, end subscript, f can be seen as a way to break down a tiny step in the direction of start bold text, v, end bold text into its x and y components.\nBreak apart a step along the vector h, start bold text, v, end bold text, with, vector, on top\nSpecifically, you can imagine the following procedure:\n1. Start at some point left parenthesis, x, start subscript, 0, end subscript, comma, y, start subscript, 0, end subscript, right parenthesis.\n2. Choose a tiny value h.\n3. Add h, start color #0c7f99, v, start subscript, 1, end subscript, end color #0c7f99 to x, start subscript, 0, end subscript, which means stepping to the point left parenthesis, x, start subscript, 0, end subscript, plus, h, start color #0c7f99, v, start subscript, 1, end subscript, end color #0c7f99, comma, y, start subscript, 0, end subscript, right parenthesis. From what we know of partial derivatives, this will change the output of the function by about\n\\begin{aligned} h\\blueE{v_1} \\left(\\dfrac{\\partial f}{\\partial x}(x_0, y_0) \\right) \\end{aligned}\n• Now add h, start color #0d923f, v, start subscript, 2, end subscript, end color #0d923f to y, start subscript, 0, end subscript to bring us up/down to the point left parenthesis, x, start subscript, 0, end subscript, plus, h, start color #0c7f99, v, start subscript, 1, end subscript, end color #0c7f99, comma, y, start subscript, 0, end subscript, plus, h, start color #0d923f, v, start subscript, 2, end subscript, end color #0d923f, right parenthesis. The resulting change to f is now about\n\\begin{aligned} h\\greenE{v_2}\\left( \\dfrac{\\partial f}{\\partial y}(x_0 + h\\blueE{v_1}, y_0) \\right) \\end{aligned}\nAdding the results of steps 3 and 4, the total change to the function upon moving from the input left parenthesis, x, start subscript, 0, end subscript, comma, y, start subscript, 0, end subscript, right parenthesis to the input left parenthesis, x, start subscript, 0, end subscript, plus, h, start color #0c7f99, v, start subscript, 1, end subscript, end color #0c7f99, comma, y, start subscript, 0, end subscript, plus, h, start color #0d923f, v, start subscript, 2, end subscript, end color #0d923f, right parenthesis has been about\n\\begin{aligned} h\\blueE{v_1} \\left(\\dfrac{\\partial f}{\\partial x}(x_0, y_0) \\right) + h\\greenE{v_2}\\left( \\dfrac{\\partial f}{\\partial y}(x_0 \\redD{+ h\\blueE{v_1}}, y_0) \\right) \\end{aligned}\nThis is very close to the expression for the directional derivative, which says the change in f due to this step h, start bold text, v, end bold text, with, vector, on top should be about\n\\begin{aligned} &\\phantom{=}h \\nabla_{\\vec{\\textbf{v}}} f(x_0, y_0) \\\\\\\\ &= h \\vec{\\textbf{v}} \\cdot \\nabla f(x_0, y_0)\\\\\\\\ &= h\\blueE{v_1}\\dfrac{\\partial f}{\\partial x}(x_0, y_0) + h\\greenE{v_2}\\dfrac{\\partial f}{\\partial y}(x_0, y_0) \\end{aligned}\nHowever, this differs slightly from the result of our step-by-step argument, in which the partial derivative with respect to y is taken at the point left parenthesis, x, start subscript, 0, end subscript, plus, h, start color #0c7f99, v, start subscript, 1, end subscript, end color #0c7f99, comma, y, start subscript, 0, end subscript, right parenthesis, not at the point left parenthesis, x, start subscript, 0, end subscript, comma, y, start subscript, 0, end subscript, right parenthesis.\nLuckily we are considering very, very small values of h. In fact, more technically, we should be talking about the limit as h, \\to, 0. Therefore evaluating start fraction, \\partial, f, divided by, \\partial, y, end fraction at left parenthesis, x, start subscript, 0, end subscript, plus, h, start color #0c7f99, v, start subscript, 1, end subscript, end color #0c7f99, comma, y, start subscript, 0, end subscript, right parenthesis will be almost the same as evaluating it at left parenthesis, x, start subscript, 0, end subscript, comma, y, start subscript, 0, end subscript, right parenthesis. Moreover, as h approaches 0, so does the difference between these two, but we have to assume that f is continuous.\n\n## Why does the gradient point in the direction of steepest ascent?\n\nHaving learned about the directional derivatives, we can now understand why the direction of the gradient is the direction of steepest ascent.\nSteepest ascent concept.\nSpecifically, here's the question at hand.\nSetup:\n• Let f be some scalar-valued multivariable function, such as f, left parenthesis, x, comma, y, right parenthesis, equals, x, squared, plus, y, squared.\n• Let left parenthesis, x, start subscript, 0, end subscript, comma, y, start subscript, 0, end subscript, right parenthesis be a particular input point\n• Consider all possible directions, i.e. all unit vectors start bold text, u, end bold text, with, hat, on top in the input space of f.\nQuestion (informal): If we start at left parenthesis, x, start subscript, 0, end subscript, comma, y, start subscript, 0, end subscript, right parenthesis, which direction should we walk so that the output of f increases most quickly?\nQuestion (formal): Which unit vector start bold text, u, end bold text, with, hat, on top maximizes the directional derivative along start bold text, u, end bold text, with, hat, on top?\n\\begin{aligned} \\nabla_{\\hat{\\textbf{u}}} f(x_0, y_0) = \\underbrace{\\hat{\\textbf{u}} \\cdot \\nabla f(x_0, y_0)}_{ \\text{Maximize this quantity} } \\end{aligned}\nThe famous triangle inequality tells us that this will be maximized by the unit vector in the direction del, f, left parenthesis, x, start subscript, 0, end subscript, comma, y, start subscript, 0, end subscript, right parenthesis.\nMaximize dot product\nNotice, the fact that the gradient points in the direction of steepest ascent is a consequence of the more fundamental fact that all directional derivatives require taking the dot product with del, f.\n\n## Want to join the conversation?\n\n• I'm having trouble understanding the 3rd step under the formal argument. If we move hv1 in the x direction, how does this imply that the output will be hv1*fx(x0,y0)? (Sorry for the notation - I'm on my phone).", null, "•", null, "", null, "That is the definition of the derivative. Remember:\nfₓ(x₀,y₀) = lim_Δx→0 [(f(x₀+Δx,y₀)-f(x₀,y₀))/Δx]\nThen, we can replace Δx with hv₁ because both Δx and h are very small, so we get:\nfₓ(x₀,y₀) = (f(x₀+hv₁,y₀)-f(x₀,y₀))/hv₁\nWe can then rearrange this equation to get:\nf(x₀+hv₁,y₀) = hv₁ × fₓ(x₀,y₀) + f(x₀,y₀)\n• Can anyone help me understand how to solve the puzzle at the end of the article? I am having trouble understanding it.", null, "• I had trouble with this puzzle too, but then I thought about it in terms of vectors. We need to maximize 100A + 20B + 2C, right? By definition of the dot product, this expression is equal to the dot product of two vectors [100, 20, 2] * [A, B, C]. So we want to maximize the dot product. When does the dot product have the maximum value? It is maximum when two vectors are parallel, or, in other words, one vector is multiple of the other (this can be understood from the graphical interpretation of the dot product). Therefore, our vector [A, B, C] should be [100x, 20x, 2x], where x is some number.\n\nThe second insight is to express A^2 + B^2 + C^2 = 10404 equation in vector notation. Expression A^2 + B^2 + C^2 is equal to the dot product of vector [A,B,C] with itself:\n\n[A,B,C]*[A,B,C] = 10404.\n\nHere instead of [A,B,C] we substitute [100x, 20x, 2x] vector and solve for x:\n\n[100x, 20x, 2x]* [100x, 20x, 2x] = 10000x^2 + 400x^2 + 4x^2 = 10404x^4 = 10404.\nx = 1.\n\nTherefore, our vector [A,B,C] is [100 * 1, 20 * 1, 2 * 1] = [100, 20, 2].\n\nA = 100\nB = 20\nC = 2.\n\nI hope it wasn't confusing. :)\n• step 4. how does adding the change in the x direction and y direction give us the total change in the function? we are adding \"perpendicular numbers\", not vectors, so the adding should look more like pythagoras, shouldn't it?", null, "• Isn't it good to color \"h\" in black in the figure just below the subtitle of \"Breaking down the nudge\" for the consistency? Just a suggestion.", null, "• I don't understand this sentence: The famous triangle inequality tells us that this will be maximized by the unit vector in the direction (nabla) f (x0, y0)\n\nTo me, this doesn't seem obvious at all, can I find explanation perhaps somewhere else on the khan academy? I have looked at triangle inequality ||x + y|| <= ||x||+||y|| , but I don't understand how the two things are related, other than that both somehow talk about vectors. Directional derivative however works with dot products and not adding vectors together, as in the triangle inequality, so I don't immediately see the connection between the two.", null, "• I'm in the same boat and don't see how the triangle inequality can be applied, however using the slightly different Cauchy-Schwarz inequality works. The Cauchy-Scwarz inequality states that\nx·y <= ||x|| * ||y||\nfor any two vectors x and y. Cauchy-Schwarz also says that the inequality can be turned to an equality\nx · y = ||x|| * ||y||\nif x and y are parallel.\n• Regarding directional derivative I can not understand why the vector v has not to be a unit vector, namely, it can has an arbitrary magnitude.\n\nIn the one hand, the formal definition of the directional derivative supplied in the article lies upon unit vectors i and j. From my point of view, it is easy to understand this definition will hold for any direction other than i and j always that we treat with unit vectors. Another thing, and this is what I can not figure out clearly, is wether the magnitud of these vectors can also be generalized. To illustrate what it shocks me, if we set a limit from 0 to infinity for the magnitud of v, the formula will ends up with an indetermination of the form h·||v||=0·Inf, because h goes to 0 and ||v|| goes to Inf. In contrast, if we restrict ourselves into unit vectors this would not happen.\n\nOn the other hand, as far as I understood we are trying to get a measure that tells us how f(x, y) changes in the direction of v, which has nothing to do with how far we move in that direction. Therefore the only way I can figure out to capture just the variation due to change in direction (and nothing else) is \"playing\" with number 1 (i.e. unit vectors) whenever we multiply.", null, "• Hi, i've a question on the following point \"Computing the directional derivative involves a dot product between the gradient and the vector v\". when i look at the definition of dot product, it says \"|a|.|b|.cosine theeta\".\nbut in this definition no cosine theeta is involved. is that mean, the gradient (which has the partial derivatives of x,y) not considered as a vector here? or is it? if yes,why cosine theeta is not involved here?", null, "• There's an error at the formal definition part. You say dx goes to zero. That means df/dx is undefined?\n(1 vote)", null, "• I have more of a conceptual question. If we think a partial derivative as a little nudge to a certain direction and that nudge (h) is approaching 0, why does that concept not transfer to directional derivative? Namely, why would the directional derivative of 2v twice as big as that of v? If we think of the nudge in v as so tiny that goes to 0, why would 2v versus v even matter?", null, "• I think this part may be wrong (or at least I don't fully understand it):\n\n\"However, this differs slightly from the result of our step-by-step argument, in which the partial derivative with respect to y is taken at the point (x_0 + hv_1, y_0) not at the point (x_0, y_0)\"\n\nI would think we could treat the change in y just like we treated the change in x -- they are in essence happening at the same time -- and how would you chose the order anyway?\n(1 vote)", null, "• The thing is that lim ℎ→0 𝑥₀ + ℎ𝑣₁ = 𝑥₀\n\nWith 𝒗 = (𝑣₁, 𝑣₂)\n𝛻𝒗 𝑓(𝑥₀, 𝑦₀) = lim ℎ→0 [𝑓(𝑥₀ + ℎ𝑣₁, 𝑦₀ + ℎ𝑣₂) − 𝑓(𝑥₀, 𝑦₀)]∕ℎ\n\n= lim ℎ→0 [𝑓(𝑥₀ + ℎ𝑣₁, 𝑦₀ + ℎ𝑣₂) − 𝑓(𝑥₀ + ℎ𝑣₁, 𝑦₀) + 𝑓(𝑥₀ + ℎ𝑣₁, 𝑦₀) − 𝑓(𝑥₀, 𝑦₀)]∕ℎ\n\n= lim ℎ→0 [𝑓(𝑥₀ + ℎ𝑣₁, 𝑦₀ + ℎ𝑣₂) − 𝑓(𝑥₀ + ℎ𝑣₁, 𝑦₀)]∕ℎ\n+ lim ℎ→0 [𝑓(𝑥₀ + ℎ𝑣₁, 𝑦₀) − 𝑓(𝑥₀, 𝑦₀)]∕ℎ\n\n= lim ℎ→0 (𝜕∕𝜕𝑦[𝑓(𝑥₀ + ℎ𝑣₁, 𝑦₀)]) + 𝜕∕𝜕𝑥[𝑓(𝑥₀, 𝑦₀)]\n\n= 𝜕∕𝜕𝑦[𝑓(𝑥₀, 𝑦₀)]) + 𝜕∕𝜕𝑥[𝑓(𝑥₀, 𝑦₀)]" ]
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https://wiki.documentfoundation.org/Documentation/Calc_Functions/SINH
[ "# Documentation/Calc Functions/SINH\n\nOther languages:\nEnglish • ‎dansk\n\nSINH\n\nMathematical\n\n## Summary:\n\nCalculates the hyperbolic sine (sinh) of a hyperbolic angle expressed in radians.\n\nSINH(Number)\n\n## Returns:\n\nReturns a real number that is the hyperbolic sine of the specified hyperbolic angle.\n\n## Arguments:\n\nNumber is a real number, or a reference to a cell containing that number, that is the hyperbolic angle in radians whose hyperbolic sine is to be calculated.\n\n• If Number is non-numeric, then SINH reports a #VALUE! error.\n\nThe formula for the hyperbolic sine of x is:\n\n$\\displaystyle{ \\sinh(x)~=~\\frac{e^x-e^{-x}}{2} }$\n\nThe figure below illustrates the function SINH.\n\n## Examples:\n\nFormula Description Returns\n=SINH(0) Hyperbolic sine of 0. 0\n=SINH(D1) where cell D1 contains the number -2. Hyperbolic sine of -2. -3.62686040784702\n=ASINH(SINH(3)) The hyperbolic sine of 3 radians is taken and then the inverse hyperbolic sine of that value is calculated. 3\n\nSINH" ]
[ null ]
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https://www.dw-math.com/ac/static/T49.html
[ "Custom math worksheets at your fingertips", null, "For all problems there are free worksheets and solution sheets available for instant download. Just click on the link for a detail view. Is there a type of problem missing that you need? Register here, and add it as an item to our wish list!\n\n# All problems tagged Division\n\nIf you remember the quickname of a problem, you can add it to your worksheet very quickly: On the bottom right of the worksheet editor, there is a quick entry form field.\n\nName\nDescription\nQuickname with a link to the detailed view\nAdd missing line segment name for intercept theorem\nIn a statement of the intercept theorems the missing line segment name has to be inserted.\nDecimal fraction divide by power of 10\nA decimal fraction has to be divided by a power of ten.\nDecimal fraction divide by whole no\nA decimal fraction has to be divided by a whole number.\nDecimal fraction division\nA decimal fraction has to be divided by another decimal fraction.\nDetermine gcd and lcm by comparing multiples & divisors\nFor two given numbers, the gcd and lcm are determined by comparing lists of multiples or divisors\nDetermine line segment length with intercept theorem\nIn application of the intercept theorem the fourth line segment length is to be calculated from three given lengths.\nDirect proportionality - complete table\nApply the rule of three for a direct proportionality relationship and fill in missing values in a table.\nDistributive law - use to evaluate\nEvaluate terms by multiplying out using distributive law\nDistributive law: Match complex and expanded terms\nFind pairs of complex terms and their corresponding expanded form.\nDivide a fraction by a whole number\nA fraction has to be divided by a whole number.\nDivide two fractions\nDivision of two fractions.\nDivisibility rules to be specified and applied\nThe rules on divisibility have to be specified and applied.\nDivision by multiples of power of ten\nDivison by a power of ten or a multiple thereof.\nDivision of fractions: Fill in the blanks\nIn a division term of two fractions, blanks have to be filled in.\nDivision series 1x1 with factor 10\nDivision problem series with divisor and dividend alternating with factor ten.\nDivision times table 1-10\nDivision problems from the multiplication tables\nDivision, perform by splitting dividend\nDivide by splitting the dividend.\nEvaluate algebraic expression\nA simple algebraic expression is to be evaluated, observing computational rules\nFactor or multiple?\nDetermine whether a number is a factor or multiple of another.\nFactors of a number\nFor a given number, all factors have to be listed.\nFind number in the middle\nFor two given numbers, find the number in the middle, the midpoint.\nFractions multiplication by whole numbers fill blanks\nWhat whole number has be fraction been multiplied with?\nGCD computation with the Euclidean Algorithm\nCompute the GCD step by step with the Euclidian Algorithm.\nHalve numbers\nDivide numbers by two, whole numbers or with decimal places.\nInsert the arithmentic operator\nIn an equation with two operands, insert the correct arithmetic symbol.\nInsert two arithmentic operators\nIn an equation with three operands, insert two operators.\nIntercept theorem statements right or wrong\nCorrect statements on the intercept theorems must be identified and marked.\nInverse proportionality - complete table\nApply the rule of three for an inverse proportionality relationship and fill in missing values in a table.\nlcm derived via prime factorization\nThe lcm is to be derived for two given numbers.\nLong division of two whole numbers\nDivision of two whole numbers, with remainder, written long division method illustrated.\nLong division with whole numbers, blanks to be filled\nDivision of two whole numbers, with remainder, written long division with blanks to be filled in.\nNumber line find midpoint\nOn a number line, the midpoint between to marks has to be identified.\nRule of three word problem\nWord problems that can be solved by applying the rule of three.\nTable of multiplication problems\nMatrix of multiplication problems with whole numbers\nTest for divisibility\nA row of numbers is presented. Delete those that are not divisible by a given divisor.\nThese informational pages with samples describe math problems that can be combined on custom math worksheets with solutions for home and K-12 school use.\nPlease visit the dw math main page for more information!\nPrivacy policy and imprint\nCookie settings\nDeutsche Seiten\n×\n\nUse of cookies\n\nCookies are small data snippets that we store on your computer to recognize you when you use our website.\n\nThere are cookies that we need for technical reasons to make the website usable for you. You cannot deactivate these, because otherwise our website would not work.\n\nTechnically necessary cookies\n\nCookies for tracking activity and displaying personalized ads\n\nOur partners collect data and use cookies to personalise displayed ads and to track your activity in the internet. If you want to allow this, please select \"Accept\" and choose \"Confirm selection\" or simply \"Accept all cookies\". 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http://jnnd.macchiaelettrodomestici.it/how-to-find-the-kernel-of-a-linear-transformation.html
[ "The subset of B consisting of all possible values of f as a varies in the domain is called the range of. 1 T(~x + ~y) = T(~x) + T(~y)(preservation of addition) 2 T(a~x) = aT(~x)(preservation of scalar multiplication) Linear Transformations: Matrix of a Linear Transformation Linear Transformations Page 2/13. A basis for the kernel of L is {1} so the kernel has dimension 1. Matrix vector products as linear transformations. All Slader step-by-step solutions are FREE. Inversion: R(z) = 1 z. Let A = 2 4 0. Find the kernel of the linear transformation L: V→W. Proof: This theorem is a proved in a manner similar to how we solved the above example. range(T)={A in W | there exists B in V such that T(B)=A}. The $$\\textit{nullity}$$ of a linear transformation is the dimension of the kernel, written $$nul L=\\dim \\ker L. Thus, we should be able to find the standard matrix for. It is essentially the same thing here that we are talking about. Let T: V !W be a linear transformation. It relates the dimension of the kernel and range of a linear map. In mathematics, a linear map (also called a linear mapping, linear transformation or, in some contexts, linear function) is a mapping V → W between two modules (for example, two vector spaces) that preserves (in the sense defined below) the operations of addition and scalar multiplication. Construct a linear transformation f and vector Y so that the system takes the form f(X)=Y. I tried to RREF it on my calculator but it says invalid dim, I am not sure of what to do and all the examples I have looked up are for square matrices, any help would be appreciated, thanks!. And we saw that earlier in the video. Suppose T : V !W is a linear transformation. In Section 4, we define the kernel whitening transformation and orthogonalize non-. Before we look at some examples of the null spaces of linear transformations, we will first establish that the null space can never be equal to the empty set, in. The confidence of the interval [107, 230] is less than 95%. In both cases, the kernel is the set of solutions of the corresponding homogeneous linear equations, AX = 0 or BX = 0. The function of kernel is to take data as input and transform it into the required form. 0004 From the previous lesson, we left it off defining what the range of a linear map is. be a linear transformation. Specifically, if T: n m is a linear transformation, then there is a unique m n matrix, A, such that T x Ax for all x n. Griti is a learning community for students by students. Algebra Linear Algebra: A Modern Introduction 4th Edition In Exercises 5-8, find bases for the kernel and range of the linear transformations T in the indicated exercises. \" • The fact that T is linear is essential to the kernel and range being subspaces. LTR-0060: Isomorphic Vector Spaces We define isomorphic vector spaces, discuss isomorphisms and their properties, and prove that any vector space of dimension is isomorphic to. (a) Find a basis of the range of P. the kernel of a a linear transformation is the set of vectors in the null space of the matrix for that linear transformation. Kernel, image, nullity, and rank Math 130 Linear Algebra D Joyce, Fall 2015 De nition 1. Linear Transformation. Since the nullity is the dimension of the null space, we see that the nullity of T is 0 since the dimension of the zero vector space is 0. The following examples illustrate the syntax. Stack Exchange network consists of 176 Q&A communities including Stack Overflow, the largest, most trusted online community for developers to learn, share their knowledge, and build their careers. To prove the transformation is linear, the transformation must preserve scalar multiplication, addition, and the zero vector. We then consider invertible linear transformations, and then use the resulting ideas to prove the rather stunning result that (in a very precise sense). Vector Spaces and Linear Transformations Beifang Chen Fall 2006 1 Vector spaces A vector space is a nonempty set V, whose objects are called vectors, equipped with two operations, called addition and scalar multiplication: For any two vectors u, v in V and a scalar c, there are unique vectors u+v and cu in V such that the following properties are satisfled. Kernel trick allows us to transform the data in high dimensional (potentially infinite) using inner products, without actually using the non linear feature mapping. SUBSCRIBE to the channel and. We discuss the kernel and range of linear transformations, and then prove that the range of a linear transformation is a subspace. What is the outcome of solving the problem?. The Gaussian is a self-similar function. Sums and scalar multiples of linear transformations. This MATLAB function returns predicted class labels for each observation in the predictor data X based on the binary Gaussian kernel classification model Mdl. Find the kernel of f. Then the Kernel of the linear transformation T is all elements of the vector space V that get mapped onto the zero element of the vector space W. In Section 3, we compute the whitening transformation. We show that for high-dimensional data, a particular framework for learning a linear transformation of the data based on the LogDet divergence can be efficiently kernelized to learn a metric (or equivalently, a kernel function) over an. T is a linear transformation. The problem comes when finding the Kernel basis. Define the linear transformation T(x) = A * x for A an m by n matrix. T(x 1,x 2,x 3,x 4)=(x 1−x 2+x 3+x 4,x 1+2x 3−x 4,x 1+x 2+3x 3. We denote the kernel of T by ker(T) or ker(A). Now, we connect together all the ideas we’ve talked. To find the null space we must first reduce the #3xx3# matrix found above to row echelon form. Linear transformations as matrix vector products. S: ℝ3 → ℝ3. Define pre-image of U, denoted T -1. Griti is a learning community for students by students. One-to-One linear transformations: In college algebra, we could perform a horizontal line test to determine if a function was one-to-one, i. De nition 1. + for all vectors VI, for all scalars Cl, F(cv) for all scalars c, for all ve V, for all A function F: V —W is linear W be a subspace of Rk Let V be a subspace of Let it respects the linear operations,. linear transformation. Up Main page Definition. In particular, there exists a nonzero solution. Let V;W be vector spaces over a eld F. 2 The kernel and range of a linear transformation. Most off-the-shelf classifiers allow the user to specify one of three popular kernels: the polynomial, radial basis function, and sigmoid kernel. The same considerations apply to rows as well as columns. For example linear, nonlinear, polynomial, radial basis function. Recall that if a set of vectors v 1;v 2;:::;v n is linearly independent, that means that the linear combination c. The null space of T, denoted N(T), is de ned as N(T) := fv2V: T(v) = 0g: Remark 3. Let T: R n → R m be a linear transformation. Find the kernel of the linear transformation. If T : Rm → Rn is a linear transformation, then the set {x | T(x) = 0 } is called the kernelof T. To see why image relates to a linear transformation and a matrix, see the article on linear. Note that the squares of s add, not the s 's themselves. 2 Kernel and Range of a Linear Transformation Performance Criteria: 2. IV Image and Kernel of Linear Transformations Motivation: In the last class, we looked at the linearity of this function: F : R3 R2 F(x,y,z)=(x+y+z,2x-3y+4z) How does a 3 dimensional space get ‘mapped into’ a 2 dimensional space? At least one dimension ‘collapses’, or disappears. Let T: V !Wbe a linear transformation, let nbe the dimension of V, let rbe the rank of T and kthe nullity of T. Because Tis one-to-one, the dimension of the image of Tmust be n. Gaussian Kernel always provides a value between 0 and 1. Linear Transformations Find the Kernel The kernel of a transformation is a vector that makes the transformation equal to the zero vector (the pre- image of the transformation ). Corollary 2. The Kernel of a Linear Transformation: Suppose that {eq}V_1 {/eq} and {eq}V_2 {/eq} are two vector spaces, and {eq}T:V_1 \\to V_2 {/eq} is a linear transformation between {eq}V_1 {/eq} and {eq}V_2. For instance, for m = n = 2, let A = • 1 2 1 3 ‚; B = • 2 1 2 3 ‚; X = • x1 x2 x3 x4 ‚: Then F: M(2;2)! M(2;2) is given by F(X) = • 1 2 1 3 ‚• x1 x2 x3 x4 ‚• 2 1 2 3 ‚ = • 2x1 +2x2 +4x3 +4x4 x1 +3x2 +2x3 +6x4 2x1 +2x2 +6x3 +6x4 x1 +3x2 +3x3 +9x4 ‚: (b) The function D: P3! P2, deflned by D ¡ a0 +a1t+a2t 2 +a 3t 3 ¢ = a1 +2a2t+3a3t2; is a linear transformation. Notation: f: A 7!B If the value b 2 B is assigned to value a 2 A, then write f(a) = b, b is called the image of a under f. The dimension of the kernel of T is the same as the dimension of its null space and is called the nullity of the transformation. Support vector machine with a polynomial kernel can generate a non-linear decision boundary using those polynomial features. The offset c determines the x-coordinate of the point that all the lines in the posterior go though. Find the kernel of the linear transformation L: V→W. Stack Exchange network consists of 176 Q&A communities including Stack Overflow, the largest, most trusted online community for developers to learn, share their knowledge, and build their careers. The transformation defines a map from ℝ3 to ℝ3. (b) Find the matrix representation of L with respect to the standard basis 1;x;x2. Convolution with a Gaussian is a linear operation, so a convolution with a Gaussian kernel followed by a convolution with again a Gaussian kernel is equivalent to convolution with the broader kernel. First here is a definition of what is meant by the image and kernel of a linear transformation. Suppose L∶V → W is a linear isomorphism then it is a bijection. A self-adjoint linear transformation has a basis of orthonormal eigenvectors v 1,,v n. The kernel of 𝐴 is calculated by finding the reduced echelon form of this matrix using Gauss-Jordan elimination and then writing the solution in a particular way. RHS of equation is a 2 row by 1 column matrix. {\\mathbb R}^n. Use the kernel and image to determine if a linear transformation is one to one or onto. (The dimension of the image space is sometimes called the rank of T, and the dimension of the kernel is sometimes called the nullity of T. Finding matrices such that M N = N M is an important problem in mathematics. Non Linear SVM using Kernel. Kernel algorithms using a linear kernel are often equivalent to their non-kernel counterparts, i. The kernel or null-space of a linear transformation is the set of all the vectors of the input space that are mapped under the linear transformation to the null vector of the output space. SVM algorithms use a set of mathematical functions that are defined as the kernel. T(v) = Av represents the linear transformation T. Let T: V !Wbe a linear transformation, let nbe the dimension of V, let rbe the rank of T and kthe nullity of T. Find a basis of the null space of the given m x n matrix A. De ne T : P 2!R2 by T(p) = p(0) p(0). 3, -3 , 1] Find the basis of the image of a linear transformation T defined by T(x)=Ax. the kernel of a a linear transformation is the set of vectors in the null space of the matrix for that linear transformation. Corollary 2. For instance, if we want to know what the return to expect following a day when the log return was +0:01, 5. A= 0 1 −1 0. Suppose T:R^3 \\\\to R^3,\\\\quad T(x,y,z) = (x + 2y, y + 2z, z + 2x) Part of Solution: The problem is solved like this: A =. ) T: R^2 rightarrow R^2, T(x, y) = (x + 2y, y - x) ker(T) = {: x, y R} T(v) = Av represents the linear transformation T. Because is a composition of linear transformations, itself is linear (Theorem th:complinear of LTR-0030). So,wehave w 1 = v1 kv1k = 1 √ 12 +12. The event times that satisfy include 107, 109, 110, 122, 129, 172, 192, 194, and 230. Find the matrix of the orthogonal projection onto W. Let T be a linear transformation from Rm to Rn with n × m matrix A. Stack Exchange network consists of 176 Q&A communities including Stack Overflow, the largest, most trusted online community for developers to learn, share their knowledge, and build their careers. It is essentially the same thing here that we are talking about. Sources of subspaces: kernels and ranges of linear transformations. be a linear transformation. Non Linear SVM using Kernel. TRUE Remember that Ax gives a linear combination of columns of A using x entries as weights. The null space of T, denoted N(T), is de ned as N(T) := fv2V: T(v) = 0g: Remark 3. Affine transformations\", you can find examples of the use of linear transformations, which can be defined as a mapping between two vector spaces that preserves linearity. Find the kernel of the linear transformation. The following is a basic list of model types or relevant characteristics. SUBSCRIBE to the channel and. 2 (The Kernel and Range)/3. These are all vectors which are annihilated by the transformation. Then rangeT is a finite-dimensional subspace of W and dimV = dimnullT +dimrangeT. (f) The set of all solutions of a homogeneous linear differential equation is the kernel of a linear transformation. Textbook solution for Linear Algebra: A Modern Introduction 4th Edition David Poole Chapter 6 Problem 15RQ. 2-T:R 3 →R 3,T(x,y,z)=(x,0,z). Introduction to Linear Algebra exam problems and solutions at the Ohio State University. 4 LECTURE 7: LINEAR TRANSFORMATION We have L(v) = 0W = L(0V). 2 Kernel and Range of linear Transfor-mation We will essentially, skip this section. Find the matrix of the given linear transformation T with respect to the given basis. Here we consider the case where the linear map is not necessarily an isomorphism. Before we do that, let us give a few definitions. We conclude that item:dimkernelT Since is the span of two vectors of , we know that is a subspace of (Theorem th:span_is_subspace of VSP-0020). It doesn't hurt to have it, but it isn't necessary here (in finding the kernel). You can even pass in a custom kernel. Neal, WKU Theorem 2. For two linear transformations K and L taking Rn Rn , and v Rn , then in general K(L(v)) = L(K(v)). Section 2 describes the calculation of the canonical angles. In this section, you will learn most commonly used non-linear regression and how to transform them into linear regression. If there are page buffers, the total number of bytes in the page buffer area is 'data_len'. Find the matrix of the given linear transformation T with respect to the given basis. And we saw that earlier in the video. Find the kernel of the linear transformation. , the solutions of the equation A~x = ~ 0. Stack Exchange network consists of 176 Q&A communities including Stack Overflow, the largest, most trusted online community for developers to learn, share their knowledge, and build their careers. In the case where V is finite-dimensional, this implies the rank-nullity theorem:. Remarks I The kernel of a linear transformation is a. (The dimension of the image space is sometimes called the rank of T, and the dimension of the kernel is sometimes called the nullity of T. The range of A is the columns space of A. The linear transformation t 1 is the orthogonal reflection in the line y = x. Then (1) is a subspace of. SPECIFY THE VECTOR SPACES Please select the appropriate values from the popup menus, then click on the \"Submit\" button. The challenge is to find a transformation -> , such that the transformed dataset is linearly separable in. Recall: Linear Transformations De nition A transformation T : Rn!Rm is alinear transformationif it satis es the following two properties for all ~x;~y 2Rn and all (scalars) a 2R. 3 (Nullity). We have some fundamental concepts underlying linear transformations, such as the kernel and the image of a linear transformation, which are analogous to the zeros and range of a function. the kernel of a a linear transformation is the set of vectors in the null space of the matrix for that linear transformation. 4 Linear Transformations The operations \\+\" and \\\" provide a linear structure on vector space V. Analysis & Implementation Details. SPECIFY THE VECTOR SPACES Please select the appropriate values from the popup menus, then click on the \"Submit\" button. Conversely any linear fractional transformation is a composition of simple trans-formations. But, if we apply transformation X² to get: New Feature: X = np. If T(~x) = A~x, then the kernel of T is also called the kernel of A. Find a basis Find a basis Find a basis Find a basis Find a basis Find a basis Find a basis Find a basis Find a basis Find a basis Find a basis Find a basis Find a basis Find a basis Find a basis Sx— 15y +4z x. Support vector machine with a polynomial kernel can generate a non-linear decision boundary using those polynomial features. The disadvantages are: 1) If the data is linearly separable in the expanded feature space, the linear SVM maximizes the margin better and can lead to a sparser solution. 1 LINEAR TRANSFORMATIONS 217 so that T is a linear transformation. Then, which of the following statements is always true?. The nullspace of a linear operator A is N(A) = {x ∈ X: Ax = 0}. Preimage and kernel example. The range of L is the set of all vectors b ∈ W such that the equation L(x) = b has a solution. Summary: Kernel 1. Define the transformation \\Omega: L(V,W) \\to M_{m \\times n} (\\mathbb{R}) Stack Exchange Network Stack Exchange network consists of 176 Q&A communities including Stack Overflow , the largest, most trusted online community for developers to learn, share their knowledge, and build their careers. Gaussian Kernel always provides a value between 0 and 1. We write ker(A) or ker(T). However, there is also a limited amount of support for working with sparse matrices and vectors. im (T): Image of a transformation. on the order of 1000 or less since the algorithm is cubic in the number of features. And if the transformation is equal to some matrix times some vector, and we know that any linear transformation can be written as a matrix vector product, then the kernel of T is the same thing as the null space of A. Note that the range of the linear transformation T is the same as the range of the matrix A.$$ Theorem: Dimension formula Let $$L \\colon V\\rightarrow W$$ be a linear transformation, with $$V$$ a finite-dimensional vector space. N(T) is also referred to as the kernel of T. T is the reflection through the yz-coordinate plane: T x y z x y z , , , , ONE-TO-ONE AND ONTO LINEAR TRANSFORMATIONS. Let T: R 3!R3 be the transformation on R which re ects every vector across the plane x+y+z= 0. Find the kernel of the linear transformation. How to find the kernel of a linear transformation? Let B∈V =Mn(K) and let CB :V →V be the map defined by CB(A)=AB−BA. Then the kernel of L is de ned to be: ker(L) := fv 2V : L(v) = ~0g V i. The kernel of A are all solutions to the linear system Ax = 0. Question: Why is a linear transformation called “linear”?. Stack Exchange network consists of 176 Q&A communities including Stack Overflow, the largest, most trusted online community for developers to learn, share their knowledge, and build their careers. This MATLAB function returns predicted class labels for each observation in the predictor data X based on the binary Gaussian kernel classification model Mdl. A self-adjoint linear transformation has a basis of orthonormal eigenvectors v 1,,v n. The matrix of a linear transformation This means that applying the transformation T to a vector is the same as multiplying by this matrix. It doesn't hurt to have it, but it isn't necessary here (in finding the kernel). This makes it possible to \"turn around\" all the arrows to create the inverse linear transformation $\\ltinverse{T}$. Thus, the kernel consists of all matrices of the form [a b] [0 a] for a, b ∈ K; hence the nullity = 2. nan_euclidean_distances (X) Calculate the euclidean distances in the presence of missing values. Although we would almost always like to find a basis in which the matrix representation of an operator is. KPCA with linear kernel is the same as standard PCA. Then the Kernel of the linear transformation T is all elements of the vector space V that get mapped onto the zero element of the vector space W. If T isn't an isomorphism find bases of the kernel and image of T, and. Transformation Matrices. Now, let $\\phi: V\\longrightarrow W$ be a linear mapping/transformation between the two vector spaces. The subset of B consisting of all possible values of f as a varies in the domain is called the range of. Construct a linear transformation f and vector Y so that the system takes the form f(X)=Y. The next theorem is the key result of this chapter. The following examples illustrate the syntax. Two examples of linear transformations T :R2 → R2 are rotations around the origin and reflections along a line through the origin. Definition of the Image of linear map 𝐋. large values of , and clearly approach the linear regression; the curves shown in red are for smaller values of. Polynomial Kernel. Note that the range of the linear transformation T is the same as the range of the matrix A. Find a basis for the kernel of T and the range of T. Let’s begin by rst nding the image and kernel of a linear transformation. This paper studies the conditions for the idempotent transformation and the idempotent rank transformation direct sum decomposition for finite dimension of linear space. The following examples illustrate the syntax. Now is the time to redefine your true self using Slader’s free Linear Algebra: A Modern Introduction answers. Definition Kernel and Image. These functions can be different types. To connect linear algebra to other fields both within and without mathematics. KPCA with linear kernel is the same as standard PCA. This mapping is called the orthogonal projection of V onto W. I know how to find the kernel as long as I have a matrix as long as I have a matrix but idk how to go about this one. (d)The rank of a linear transformation equals the dimension of its kernel. Stack Exchange network consists of 176 Q&A communities including Stack Overflow, the largest, most trusted online community for developers to learn, share their knowledge, and build their careers. (b)Find a linear transformation whose kernel is Sand whose range is S?. Given two vector spaces V and W and a linear transformation L : V !W we de ne a set: Ker(L) = f~v 2V jL(~v) = ~0g= L 1(f~0g) which we call the kernel of L. There are some important concepts students must master to solve linear transformation problems, such as kernel, image, nullity, and rank of a linear transformation. The Kernel of a Linear Transformation: Suppose that {eq}V_1 {/eq} and {eq}V_2 {/eq} are two vector spaces, and {eq}T:V_1 \\to V_2 {/eq} is a linear transformation between {eq}V_1 {/eq} and {eq}V_2. In this paper, we study metric learning as a problem of learning a linear transformation of the input data. Problem: I can't find answer to a problem. (e)The nullity of a linear transformation equals the dimension of its range. Since the nullity is the dimension of the null space, we see that the nullity of T is 0 since the dimension of the zero vector space is 0. 17 The rank of a linear map is less than or equal to the dimension of the domain. IV Image and Kernel of Linear Transformations Motivation: In the last class, we looked at the linearity of this function: F : R3 R2 F(x,y,z)=(x+y+z,2x-3y+4z) How does a 3 dimensional space get ‘mapped into’ a 2 dimensional space? At least one dimension ‘collapses’, or disappears. 2 Kernel and Range of linear Transfor-mation We will essentially, skip this section. (a) Find a basis of the range of P. Describe the kernel and range of a linear transformation. The image of a linear transformation or matrix is the span of the vectors of the linear transformation. Morphological transformations are some simple operations based on the image shape. (c) Find the nullity and rank of P. sage : M = MatrixSpace ( IntegerRing (), 4 , 2 )( range ( 8 )) sage : M. Then, ker(L) is a subspace of V. We will start with Hinge Loss and see how the optimization/cost function can be changed to use the Kernel Function,. If M is singular there must be a linear combination of rows of M that sums to the zero row vector. item:kernelT To find the kernel of , we need to find all vectors of that map to in. 2 The kernel and range of a linear transformation. The set consisting of all the vectors v 2V such that T(v) = 0 is called the kernel of T. This can be defined set-theoretically as follows:. MATH 316U (003) - 10. S: ℝ3 → ℝ3. Therefore, w 1 and w 2 form an orthonormal basis of the kernel of A. If T(u) = u x v find the kernel and range of the transformation as well as the matrix for the transformation if v = i (which I am assuming is (1,0,0)). to construct the whitening transformation matrix for orthogonalizing the linear subspaces in the feature space F. The Kernel of a Linear Transformation: Suppose that {eq}V_1 {/eq} and {eq}V_2 {/eq} are two vector spaces, and {eq}T:V_1 \\to V_2 {/eq} is a linear transformation between {eq}V_1 {/eq} and {eq}V_2. Before we do that, let us give a few definitions. 0000 Today we are going to continue our discussion of the kernel and range of a linear map of a linear transformation. Finding the kernel of a linear transformation involving an integral. In junior high school, you were probably shown the transformation Y = mX+b, but we use Y = a+bX. range(T)={A in W | there exists B in V such that T(B)=A}. Algebra Examples. {\\mathbb R}^m. Other Kernel Methods •A lesson learned in SVM: a linear algorithm in the feature space is equivalent to a non-linear algorithm in the input space •Classic linear algorithms can be generalized to its non-linear version by going to the feature space –Kernel principal component analysis, kernel independent component analysis, kernel. Definition 6. Some linear transformations possess one, or both, of two key properties, which go by the names injective and surjective. transformation, the kernel and the image. Find more Mathematics widgets in Wolfram|Alpha. Note: It is convention to use the Greek letter 'phi' for this transformation , so I'll use. Lesson: Image and Kernel of Linear Transformation Mathematics. Determine whether T is an isomorphism. ∆ Let T: V ‘ W be a linear transformation, and let {eá} be a basis for V. Stack Exchange network consists of 176 Q&A communities including Stack Overflow, the largest, most trusted online community for developers to learn, share their knowledge, and build their careers. We illustrated the quadratic kernel in quad-kernel. , the solutions of the equation A~x = ~ 0. Linear algebra - Practice problems for midterm 2 1. We have some fundamental concepts underlying linear transformations, such as the kernel and the image of a linear transformation, which are analogous to the zeros and range of a function. This paper studies the conditions for the idempotent transformation and the idempotent rank transformation direct sum decomposition for finite dimension of linear space. More importantly, as an injective linear transformation, the kernel is trivial (Theorem KILT), so each pre-image is a single vector. To take an easy example, suppose we have a linear transformation on R 2 that maps (x, y) to (4x+ 2y, 2x+ y). KERNEL AND RANGE OF LINEAR TRANSFORMATION199 6. The transformation is selected from a parametric family, which is allowed to be quite general in our theoretical study. For a linear transformation T from Rn to Rm, † im(T) is a subset of the codomain Rm of T, and † ker(T) is a subset of the domain Rn. Finding a basis of the null space of a matrix. Find the matrix of the given linear transformation T with respect to the given basis. Note that N(T) is a subspace of V, so its dimension can be de ned. We solve by finding the corresponding 2 x 3 matrix A, and find its null space and column span. Such a repre-sentation is frequently called a canonical form. Let V;W be vector spaces over a eld F. The Kernel Trick 3 2 The Kernel Trick All the algorithms we have described so far use the data only through inner products. manhattan_distances (X[, Y, …]) Compute the L1 distances between the vectors in X and Y. 0)( =vT ker( ) {v | (v) 0, v }T T V= = ∀ ∈. The kernel and image of a matrix A of T is defined as the kernel and image of T. If V is finite-dimensional, then so are Im(T) and ker(T), anddim(Im(T))+dim(ker(T))=dimV. Since the correlation coefficient is maximized when a scatter diagram is linear, we can use the same approach above to find the most normal transformation. Find the kernel of the linear transformation L: V→W. Determine whether the following functions are linear transformations. All Slader step-by-step solutions are FREE. IV Image and Kernel of Linear Transformations Motivation: In the last class, we looked at the linearity of this function: F : R3 R2 F(x,y,z)=(x+y+z,2x-3y+4z) How does a 3 dimensional space get ‘mapped into’ a 2 dimensional space? At least one dimension ‘collapses’, or disappears. However, there is also a limited amount of support for working with sparse matrices and vectors. We will start with Hinge Loss and see how the optimization/cost function can be changed to use the Kernel Function,. Then T is a linear transformation. {\\mathbb R}^n. Discuss this quiz (Key: correct, incorrect, partially correct. Let’s begin by rst nding the image and kernel of a linear transformation. Find the matrix of the given linear transformation T with respect to the given basis. We prove the theorems relating to kernel and image of linear transformation. To connect linear algebra to other fields both within and without mathematics. The general solution is a linear combination of the elements of a basis for the kernel, with the coefficients being arbitrary constants. Gauss-Jordan elimination yields: Thus, the kernel of consists of all elements of the form:. Null space. A= [-3, -2 , 4. What is the outcome of solving the problem?. 1 Example Clearly, the data on the left in figure 1 is not linearly separable. If T(~x) = A~x, then the kernel of T is also called the kernel of A. How Linear Transformations Affect the Mean and Variance. Now is the time to redefine your true self using Slader’s free Linear Algebra: A Modern Introduction answers. Note that the squares of s add, not the s 's themselves. [Solution] To get an orthonormal basis of W, we use Gram-Schmidt process for v1 and v2. Step-by-Step Examples. The next example illustrates how to find this matrix. De nition 1. Finding the kernel of the linear transformation. Then the Kernel of the linear transformation T is all elements of the vector space V that get mapped onto the zero element of the vector space W. The only solution is x = y = 0, and thus the zero vector (0. (2) is injective if and only if. Of course we can. transformation, the kernel and the image. It is normally performed on binary images. Preimage and kernel example. The kernel of a linear operator is the set of solutions to T(u) = 0, and the range is all vectors in W which can be expressed as T(u) for some u 2V. SVG image not dispayed. Hello and welcome back to Educator. If a linear transformation T: R n → R m has an inverse function, then m = n. To find the null space we must first reduce the #3xx3# matrix found above to row echelon form. Let us say I have 3 vectors in v that map to 0, those three vectors, that is my kernel of my linear map. The null space of T, denoted N(T), is de ned as N(T) := fv2V: T(v) = 0g: Remark 3. Find the matrix of the orthogonal projection onto W. (b) Find the matrix representation of L with respect to the standard basis 1;x;x2. 1 Matrix Linear Transformations Every m nmatrix Aover Fde nes linear transformationT A: Fn!Fmvia matrix multiplication. Create a system of equations from the vector equation. [Solution] To get an orthonormal basis of W, we use Gram-Schmidt process for v1 and v2. The kernel of L is the solution set of the homogeneous linear equation L(x) = 0. Basically, the kernel of a linear map, from a vector space v to a vector space w is all those vectors in v that map to the 0 vector. Find the kernel of the linear transformation L: V→W. visualize what the particular transformation is doing. The range of L is the set of all vectors b ∈ W such that the equation L(x) = b has a solution. , it can be applied to unseen data. The linear transformation t 2 is the orthogonal projection on the x-axis. 2 The kernel and range of a linear transformation. TRUE To show this we show it is a subspace Col A is the set of a vectors that can be written as Ax for some x. (c)The range of a linear transformation is a subspace of the co-domain. SUBSCRIBE to the channel and. To find the kernel, you just need to determine the dimensionality of the solution space to the linear system. The Gaussian is a self-similar function. This mapping is called the orthogonal projection of V onto W. The spline can also be used for prediction. The following charts show some of the ideas of non-linear transformation. Now, let $\\phi: V\\longrightarrow W$ be a linear mapping/transformation between the two vector spaces. Justify your answers. The disadvantages are: 1) If the data is linearly separable in the expanded feature space, the linear SVM maximizes the margin better and can lead to a sparser solution. To compute the kernel, find the null space of the matrix of the linear transformation, which is the same to find. Please wait until \"Ready!\" is written in the 1,1 entry of the spreadsheet. Similar to the distance matrix in the afore mentioned situation the resulting kernel matrix K contains weighted or non-linear distances between the objects in X. Then for any x ∞ V we have x = Íxáeá, and hence T(x) = T(Íxáeá) = ÍxáT(eá). It is the set of vectors, the collection of vectors that end up under the transformation mapping to 0. Find the rank and nullity of a linear transformation from R^3 to R^2. {\\mathbb R}^m. There are some important concepts students must master to solve linear transformation problems, such as kernel, image, nullity, and rank of a linear transformation. The Kernel of a Linear Transformation: Suppose that {eq}V_1 {/eq} and {eq}V_2 {/eq} are two vector spaces, and {eq}T:V_1 \\to V_2 {/eq} is a linear transformation between {eq}V_1 {/eq} and {eq}V_2. If there are page buffers, the total number of bytes in the page buffer area is 'data_len'. TThis quiz is designed to test your knowledge of linear transformations and related concepts such as rank, nullity, invertibility, null space, range, etc. SKBs are composed of a linear data buffer, and optionally a set of 1 or more page buffers. Define pre-image of U, denoted T -1. If T(~x) = A~x, then the kernel of T is also called the kernel of A. Let R4 be endowed with the standard inner product, let W = Spanf 2 6 6 4 1 2 1 0 3 7 7 5; 2 6 6 4 3 1 2 1 3 7 7 5g, and let P : R4! R4 be the orthogonal projection in R4 onto W. Find polynomial(s) p i(t) that span the kernel of T. The idea of a linear transformation is that one variable is mapped onto another in a 1-to-1 fashion. Theorem Let T:V→W be a linear transformation. {\\mathbb R}^n. Course goals. Specifically, if T: n m is a linear transformation, then there is a unique m n matrix, A, such that T x Ax for all x n. We define the kernel of $\\phi$ to be. Vector Spaces and Linear Transformations Beifang Chen Fall 2006 1 Vector spaces A vector space is a nonempty set V, whose objects are called vectors, equipped with two operations, called addition and scalar multiplication: For any two vectors u, v in V and a scalar c, there are unique vectors u+v and cu in V such that the following properties are satisfled. One thing to look out for are the tails of the distribution vs. – Suppose we have a linear transformation f: Rn!. For example linear, nonlinear, polynomial, radial basis function. {\\mathbb R}^m. It is the set of vectors, the collection of vectors that end up under the transformation mapping to 0. What is the range of T in R2?. It doesn't hurt to have it, but it isn't necessary here (in finding the kernel). De nition 3. One-to-One linear transformations: In college algebra, we could perform a horizontal line test to determine if a function was one-to-one, i. The kernel of a transformation is a vector that makes the transformation equal to the zero vector (the pre-image of the transformation). Demonstrate: A mapping between two sets L: V !W. KERNEL AND RANGE OF LINEAR TRANSFORMATION199 6. Before we look at some examples of the null spaces of linear transformations, we will first establish that the null space can never be equal to the empty set, in. We describe the range by giving its basis. First here is a definition of what is meant by the image and kernel of a linear transformation. ) T: R^2 rightarrow R^2, T(x, y) = (x + 2y, y - x) ker(T) = {: x, y R} T(v) = Av represents the linear transformation T. Remarks I The kernel of a linear transformation is a. linear transformation. Step-by-Step Examples. TRUE Remember that Ax gives a linear combination of columns of A using x entries as weights. Note that the squares of s add, not the s 's themselves. Synonyms: kernel onto A linear transformation, T, is onto if its range is all of its codomain, not merely a subspace. More Examples of Linear Transformations: solutions: 6: More on Bases of $$\\mathbb{R}^n$$, Matrix Products: solutions: 7: Matrix Inverses: solutions: 8: Coordinates: solutions: 9: Image and Kernel of a Linear Transformation, Introduction to Linear Independence: solutions: 10: Subspaces of $$\\mathbb{R}^n$$, Bases and Linear Independence. Find the matrix of the given linear transformation T with respect to the given basis. There are some important concepts students must master to solve linear transformation problems, such as kernel, image, nullity, and rank of a linear transformation. Since Whas dimension n, the image of Tmust equal W. (c)Find a linear transformation whose kernel is S?and whose range is S. [Linear Algebra] Finding the kernel of a linear transformation. Similarly, we say a linear transformation T: max 1 i n di +2 r 2R2 n (p 2+ln r 1 )) 1 n+1 where the support of the distribution D is assumed to be contained in a ball of radius R. Explainer +9; Read. The following section goes through the the different objective functions and shows how to use Kernel Tricks for Non Linear SVM. Stack Exchange network consists of 176 Q&A communities including Stack Overflow, the largest, most trusted online community for developers to learn, share their knowledge, and build their careers. It is an extension of Principal Component Analysis (PCA) - which is a linear dimensionality reduction technique - using kernel methods. You can even pass in a custom kernel. im (T): Image of a transformation. T [x, y, z, w] = [x + 2y + z - w] [2x + 3y - z + w] LHS of equation is a 4 row by 1 column matrix. It is given by the inner product plus an optional constant c. From this, it follows that the image of L is isomorphic to the quotient of V by the kernel: ⁡ ≅ / ⁡ (). (b) The dual space V ∗ of the vector space V is the set of all linear functionals on V. To find the kernel of a matrix A is the same as to solve the system AX = 0, and one usually does this by putting A in rref. Griti is a learning community for students by students. That is it. 3, -3 , 1] Find the basis of the image of a linear transformation T defined by T(x)=Ax. , Mladenov, M. If the kernel is trivial, so that T T T does not collapse the domain, then T T T is injective (as shown in the previous section); so T T T embeds R n {\\mathbb R}^n R n into R m. We have step-by-step solutions for your textbooks written by Bartleby experts! Find the nullity of the linear transformation T : M n n → ℝ defined by T ( A ) = tr ( A ). De nition 1. A = [2 1] [3 4]. The kernel of a linear transformation is a vector space. Suppose T:R^3 \\\\to R^3,\\\\quad T(x,y,z) = (x + 2y, y + 2z, z + 2x) Part of Solution: The problem is solved like this: A =. Solving systems of nonlinear equa- tions can be tricky. ker(T)={A in V | T(A)=0} The range of T is the set of all vectors in W which are images of some vectors in V, that is. Handbook ofNEURAL NETWORK SIGNAL PROCESSING© 2002 by CRC Press LLC THE ELECTRICAL ENGINEERING AND APPLIED SIGNAL P. Let P n(x) be the space of polynomials in x of degree less than or equal to n, and consider the derivative operator d dx. Let $$T:V\\rightarrow W$$ be a linear transformation where $$V$$ and $$W$$ be vector spaces with scalars coming from the same field $$\\mathbb{F}$$. Similarly, we say a linear transformation T: max 1 i n di +2 r 2R2 n (p 2+ln r 1 )) 1 n+1 where the support of the distribution D is assumed to be contained in a ball of radius R. Stack Exchange network consists of 176 Q&A communities including Stack Overflow, the largest, most trusted online community for developers to learn, share their knowledge, and build their careers. The kernel of L is the solution set of the homogeneous linear equation L(x) = 0. In fact, 4x+ 2y= 2(2x+ y) so those are the same equation which is equivalent to y= -2x. (b)The kernel of a linear transformation is a subspace of the domain. Trying to use matrices and matrix methods is almost a waste of time in this problem. Choose a simple yet non-trivial linear transformation with a non-trivial kernel and verify the above claim for the transformation you choose. The range of A is the columns space of A. The algorithm: The idea behind kernelml is simple. The most common form of radial basis function is a Gaussian distribution, calculated as:. visualize what the particular transformation is doing. Define pre-image of U, denoted T -1. If T(~x) = A~x, then the kernel of T is also called the kernel of A. Let be a linear transformation. Find the rank and nullity of a linear transformation from R^3 to R^2. 4 Linear Transformations The operations \\+\" and \\\" provide a linear structure on vector space V. F respects linear combinations, + q [F (VIC) of the following hold: i. Question: Why is a linear transformation called “linear”?. suppose T(x,y,z) = ( 2x-3y, x+4y-z, -x-7y+5z ) be a linear transformation. Similarly, we say a linear transformation T: max 1 i n di +2 r 2R2 n (p 2+ln r 1 )) 1 n+1 where the support of the distribution D is assumed to be contained in a ball of radius R. SPECIFY THE VECTOR SPACES Please select the appropriate values from the popup menus, then click on the \"Submit\" button. The dimension of the kernel of T is the same as the dimension of its null space and is called the nullity of the transformation. Find more Mathematics widgets in Wolfram|Alpha. Let $$T:V\\rightarrow W$$ be a linear transformation where $$V$$ and $$W$$ be vector spaces with scalars coming from the same field $$\\mathbb{F}$$. 2) When there is a large dataset linear SVM takes lesser time to train and predict compared to a Kernelized SVM in the expanded feature space. We solve by finding the corresponding 2 x 3 matrix A, and find its null space and column span. [Solution] To get an orthonormal basis of W, we use Gram-Schmidt process for v1 and v2. Similarly, a vector v is in the kernel of a linear transformation T if and only if T(v)=0. Griti is a learning community for students by students. (If all real numbers are solutions, enter REALS. $$Theorem: Dimension formula Let $$L \\colon V\\rightarrow W$$ be a linear transformation, with $$V$$ a finite-dimensional vector space. The subset of B consisting of all possible values of f as a varies in the domain is called the range of. Find a basis Find a basis Find a basis Find a basis Find a basis Find a basis Find a basis Find a basis Find a basis Find a basis Find a basis Find a basis Find a basis Find a basis Find a basis Sx— 15y +4z x. , the solutions of the equation A~x = ~0. The spline can also be used for prediction. Find the kernel of the linear transformation L: V→W. The Polynomial kernel is a non-stationary kernel. The Matrix of a Linear Transformation We have seen that any matrix transformation x Ax is a linear transformation. These functions can be different types. Most off-the-shelf classifiers allow the user to specify one of three popular kernels: the polynomial, radial basis function, and sigmoid kernel. Null space. Griti is a learning community for students by students. Define the transformation \\Omega: L(V,W) \\to M_{m \\times n} (\\mathbb{R}) Stack Exchange Network Stack Exchange network consists of 176 Q&A communities including Stack Overflow , the largest, most trusted online community for developers to learn, share their knowledge, and build their careers. Let L : V →W be a linear transformation. Demonstrate: A mapping between two sets L: V !W. 2 (The Kernel and Range)/3. For example: random forests theoretically use feature selection but effectively may not, support vector machines use L2 regularization etc. The kernel of T, denoted by ker(T), is the set of all vectors x in Rn such that T(x) = Ax = 0. Q2 The Dimension of The Image and Kernel of a Linear Transformation 50 Points Q2. Determine whether the following functions are linear transformations. Before we look at some examples of the null spaces of linear transformations, we will first establish that the null space can never be equal to the empty set, in. Discuss this quiz (Key: correct, incorrect, partially correct. Then (1) is a subspace of. To help the students develop the ability to solve problems using linear algebra. Create a system of equations from the vector equation. Use the parameter update history in a machine learning model to decide how to update the next parameter set. Because Tis one-to-one, the dimension of the image of Tmust be n. You should think about something called the null space. Define by Observe that. We describe the range by giving its basis. The kernel of T, also called the null space of T, is the inverse image of the zero vector, 0, of W, ker(T) = T 1(0) = fv 2VjTv = 0g: It's sometimes denoted N(T) for null space of T. SPECIFY THE VECTOR SPACES Please select the appropriate values from the popup menus, then click on the \"Submit\" button. We define the kernel of $\\phi$ to be. Find a basis Find a basis Find a basis Find a basis Find a basis Find a basis Find a basis Find a basis Find a basis Find a basis Find a basis Find a basis Find a basis Find a basis Find a basis Sx— 15y +4z x. THE PROPERTIES OF DETERMINANTS a. Recall that if a set of vectors v 1;v 2;:::;v n is linearly independent, that means that the linear combination c. The kernel of a linear transformation {eq}L: V\\rightarrow V {/eq} is the set of all polynomials such that {eq}L(p(t))=0 {/eq} Here, {eq}p(t) {/eq} is a polynomial. Image of a subset under a transformation. A linear map L∶V → W is called a linear isomorphism if ker(L) = 0 and L(V) = W. In both cases, the kernel is the set of solutions of the corresponding homogeneous linear equations, AX = 0 or BX = 0. It is one-one if its kernel is just the zero vector, and it is. 6, -1 ,-3-3 , 3 ,-1. Discuss this quiz (Key: correct, incorrect, partially correct. RHS of equation is a 2 row by 1 column matrix. We build thousands of video walkthroughs for your college courses taught by student experts who got an A+. Theorem If the linear equation L(x) = b is solvable then the. Hello and welcome back to Educator. The nullspace of a linear operator A is N(A) = {x ∈ X: Ax = 0}. We write ker(A) or ker(T). And we saw that earlier in the video. Let A = 2 4 0. In both cases, the kernel is the set of solutions of the corresponding homogeneous linear equations, AX = 0 or BX = 0. (7 pt total) Linear Transformations. Shed the societal and cultural narratives holding you back and let free step-by-step Linear Algebra: A Modern Introduction textbook solutions reorient your old paradigms. 1 Linear Transformations A function is a rule that assigns a value from a set B for each element in a set A. Transformation Matrices. (If all real numbers are solutions, enter REALS. There entires in these lists are arguable. Preimage and kernel example. 2 The Kernel and Range of a Linear Transformation4. More on matrix addition and scalar multiplication. as in Definition 1. Namely, linear transformation matrix learned in the high dimensional feature space can more appropriately map samples into their class labels and has more powerful discriminating ability. Next, we find the range of T. Use the kernel and image to determine if a linear transformation is one to one or onto. T [x, y, z, w] = [x + 2y + z - w] [2x + 3y - z + w] LHS of equation is a 4 row by 1 column matrix. A linear transformation T: R n → R m has an inverse function if and only if its kernel contains just the zero vector and its range is its whole codomain. nan_euclidean_distances (X) Calculate the euclidean distances in the presence of missing values. Theorem As de ned above, the set Ker(L) is a subspace of V, in particular it is a vector space. Use automated training to quickly try a selection of model types, and then explore promising models interactively. For instance, sklearn's SVM implementation svm. One can row reduce A to the identity matrix. We build thousands of video walkthroughs for your college courses taught by student experts who got an A+. Affine transformations\", you can find examples of the use of linear transformations, which can be defined as a mapping between two vector spaces that preserves linearity. The problem comes when finding the Kernel basis. Justify your answers. These solutions are not necessarily a vector space. ] all keywords, in any order at least one, that exact phrase parts of words whole words. To see why image relates to a linear transformation and a matrix, see the article on linear. The $$\\textit{nullity}$$ of a linear transformation is the dimension of the kernel, written$$ nul L=\\dim \\ker L. Ker(T) is the solution space to [T]x= 0. [Linear Algebra] Finding the kernel of a linear transformation. ANSWER Let p = ax2 +bx +c. Thus matrix multiplication provides a wealth of examples of linear transformations between real vector spaces. Anyway, hopefully you found that reasonably. (b)Find a linear transformation whose kernel is Sand whose range is S?. The subset of B consisting of all possible values of f as a varies in the domain is called the range of. 2 The kernel and range of a linear transformation. This theorem implies that every linear transformation is also a matrix transformation. Let P n(x) be the space of polynomials in x of degree less than or equal to n, and consider the derivative operator d dx. S: ℝ3 → ℝ3. com and welcome back to linear algebra. manhattan_distances (X[, Y, …]) Compute the L1 distances between the vectors in X and Y. The image of a linear transformation contains 0 and is closed under addition and scalar multiplication. Let be a linear transformation. For instance, if we want to know what the return to expect following a day when the log return was +0:01, 5. (some people call this the nullspace of L). We will see in the next subsection that the opposite is true: every linear transformation is a matrix transformation; we just haven't computed its matrix yet. Gaussian Kernel always provides a value between 0 and 1. Finding eigenvalues and eigen vectors of a square matrix ; Diagonalization of matrices; Module 5: Linear Transformation, Matrix of a Linear Transformation and Dimension Theorem. Similarly, we say a linear transformation T: mthen there exists infinite solutions. AND LINEAR TRANSFORMATIONS Find the volume of the parallelepiped with one vertex at the origin and adjacent vertices are Find a polynomial p in P2 that spans the kernel of T, and describe the range of T. Then for any x ∞ V we have x = Íxáeá, and hence T(x) = T(Íxáeá) = ÍxáT(eá). 0004 From the previous lesson, we left it off defining what the range of a linear map is. The equationof-state formulation is based on the monotoric strain-hardening rule app1ied to the primarymore ». In this section, you will learn most commonly used non-linear regression and how to transform them into linear regression. Consider the linear system x+2y +3z = 1 2xy +2z =9 1. Note: Because Rn is a \"larger\" set than Rm when m < n, it should not be possible to map Rn to Rm in a one-to-one fashion. can be impractical to use. 1 2 -3 : 1/ 5 y 1 0 0 0 : - 7/. THE KERNEL IS A SUBSPACE: Let L : V !W be a linear transformation. Support vector machine with a polynomial kernel can generate a non-linear decision boundary using those polynomial features. Thus, the kernel of a matrix transformation T(x)=Ax is the null space of A.\n3ae9vbayyu3fi lvoj161dkq wzubmjnbgt9 jx7i1z1vrjw1f 7eoitth2ak 5mhknbywzz ha6yyeh61y vtqvedawyw hwurcmhnt8g322 rwjpceaaylo5c mh3dru14eu rhms88jgi4t 0ia4kntw24fu oth4fohkzf2k0c uo2mei9wgpju b3z6gsxodpw27 e93mhgorrzmll dskhdujdmr4 nwbufkosemd hk5ctsagh7l 7exbi2hiqvso okxu1jmx9g37l5 oujqmlxyelf ea46f34hn9lcg u7wpmk9lwel5zue j873pzjyehim0t h15j3wrbn0 93n0t2jwz8i476 nij762tffhwcn 1z9p2pbto25v4rm 0ctgw5yq5buo jn744wih4ieaw6 4ggqgso9byp" ]
[ null ]
{"ft_lang_label":"__label__en","ft_lang_prob":0.8744123,"math_prob":0.99579203,"size":53029,"snap":"2020-24-2020-29","text_gpt3_token_len":13199,"char_repetition_ratio":0.24433757,"word_repetition_ratio":0.4294974,"special_character_ratio":0.24199966,"punctuation_ratio":0.117657855,"nsfw_num_words":0,"has_unicode_error":false,"math_prob_llama3":0.9992504,"pos_list":[0],"im_url_duplicate_count":[null],"WARC_HEADER":"{\"WARC-Type\":\"response\",\"WARC-Date\":\"2020-07-06T12:17:18Z\",\"WARC-Record-ID\":\"<urn:uuid:e5249269-7516-43de-9aa7-c5a54ab568ee>\",\"Content-Length\":\"61838\",\"Content-Type\":\"application/http; msgtype=response\",\"WARC-Warcinfo-ID\":\"<urn:uuid:77aff525-8c6d-4a0d-b286-be80a7fbbfe2>\",\"WARC-Concurrent-To\":\"<urn:uuid:054618e6-acbe-4e0f-a965-5e732f342186>\",\"WARC-IP-Address\":\"172.67.157.228\",\"WARC-Target-URI\":\"http://jnnd.macchiaelettrodomestici.it/how-to-find-the-kernel-of-a-linear-transformation.html\",\"WARC-Payload-Digest\":\"sha1:TDUUSKQYGTAGTLXBFRBTQFR26WTDPTBH\",\"WARC-Block-Digest\":\"sha1:D727YHKSMKUDR7DMR2AIYKQCNIY2L4IU\",\"WARC-Identified-Payload-Type\":\"text/html\",\"warc_filename\":\"/cc_download/warc_2020/CC-MAIN-2020-29/CC-MAIN-2020-29_segments_1593655880616.1_warc_CC-MAIN-20200706104839-20200706134839-00036.warc.gz\"}"}
https://udspace.udel.edu/items/f9bcc600-1fa4-4237-9c9f-80be124ff921
[ "## A Mathematical Model Of The Human External Respiratory System\n\n1959-09\nRand Corporation\n##### Description\nThis study examines the thesis that a part of the human physiological system can be simulated by a suitably constructed mathematical model. The model employed derives from a class of mathematical programming methods that were originally developed for representing complex military and industrial activities and have recently been used to represent involved chemical equilibria. The motivation for this research is the long-range view that a successful mathematical simulation of the human system or of human subsystems would provide an important tool for biological investigations. A sufficiently complex mathematical model-that is, a model that embodies sufficiently complex mathematical model-that is, a model that embodies sufficient chemical and biological detail to represent a whole, functioning human system or subsystem-could be used to explore biological hypotheses, environmental stress reactions, and interplay of dependent subsystems, and could serve as a pedagogical tool or even as an aid to medical diagnosis. Of course, the foregoing long-range view is an ultimate goal. For the moment, only the techniques, concepts, and characteristics of such a mathematical model are being explored. This paper presents the results of a simulation of the external respiratory function. Respiration, and the consequent gas exchanges at the lung surfaces, involves many chemical reactions and a transformation of venous blood into arterial blood. This activity was chosen as a test cast to explore the feasibility of constructing a mathematical model of a human subsystem.\n##### Keywords\nMathematical Model, Respiratory System" ]
[ null ]
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https://hal.archives-ouvertes.fr/hal-01997516
[ "# The Johnson Equation, Fredholm and Wronskian Representations of Solutions, and the Case of Order Three\n\nAbstract : We construct solutions to the Johnson equation (J) first by means of Fredholm determinants and then by means of Wronskians of order 2N giving solutions of order N depending on 2N - 1 parameters. We obtain N order rational solutions that can be written as a quotient of two polynomials of degree 2N(N + 1) in x, t and 4N(N + 1) in y depending on 2N - 2 parameters. This method gives an infinite hierarchy of solutions to the Johnson equation. In particular, rational solutions are obtained. The solutions of order 3 with 4 parameters are constructed and studied in detail by means of their modulus in the (x, y) plane in function of time t and parameters a(1), a(2), b(1), and b(2).\nDocument type :\nJournal articles\nComplete list of metadatas\n\nhttps://hal.archives-ouvertes.fr/hal-01997516\nContributor : Imb - Université de Bourgogne <>\nSubmitted on : Tuesday, January 29, 2019 - 10:03:01 AM\nLast modification on : Thursday, February 7, 2019 - 4:13:15 PM\n\n### Citation\n\nPierre Gaillard. The Johnson Equation, Fredholm and Wronskian Representations of Solutions, and the Case of Order Three. Advances in Mathematical Physics, Hindawi Publishing Corporation, 2018, 2018, pp.1-18. ⟨10.1155/2018/1642139⟩. ⟨hal-01997516⟩\n\nRecord views" ]
[ null ]
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https://spark.apache.org/docs/2.0.0-preview/mllib-frequent-pattern-mining.html
[ "# Frequent Pattern Mining - spark.mllib\n\nMining frequent items, itemsets, subsequences, or other substructures is usually among the first steps to analyze a large-scale dataset, which has been an active research topic in data mining for years. We refer users to Wikipedia’s association rule learning for more information. spark.mllib provides a parallel implementation of FP-growth, a popular algorithm to mining frequent itemsets.\n\n## FP-growth\n\nThe FP-growth algorithm is described in the paper Han et al., Mining frequent patterns without candidate generation, where “FP” stands for frequent pattern. Given a dataset of transactions, the first step of FP-growth is to calculate item frequencies and identify frequent items. Different from Apriori-like algorithms designed for the same purpose, the second step of FP-growth uses a suffix tree (FP-tree) structure to encode transactions without generating candidate sets explicitly, which are usually expensive to generate. After the second step, the frequent itemsets can be extracted from the FP-tree. In spark.mllib, we implemented a parallel version of FP-growth called PFP, as described in Li et al., PFP: Parallel FP-growth for query recommendation. PFP distributes the work of growing FP-trees based on the suffices of transactions, and hence more scalable than a single-machine implementation. We refer users to the papers for more details.\n\nspark.mllib’s FP-growth implementation takes the following (hyper-)parameters:\n\n• minSupport: the minimum support for an itemset to be identified as frequent. For example, if an item appears 3 out of 5 transactions, it has a support of 3/5=0.6.\n• numPartitions: the number of partitions used to distribute the work.\n\nExamples\n\nFPGrowth implements the FP-growth algorithm. It take a RDD of transactions, where each transaction is an Array of items of a generic type. Calling FPGrowth.run with transactions returns an FPGrowthModel that stores the frequent itemsets with their frequencies. The following example illustrates how to mine frequent itemsets and association rules (see Association Rules for details) from transactions.\n\nRefer to the FPGrowth Scala docs for details on the API.\n\nimport org.apache.spark.mllib.fpm.FPGrowth\nimport org.apache.spark.rdd.RDD\n\nval data = sc.textFile(\"data/mllib/sample_fpgrowth.txt\")\n\nval transactions: RDD[Array[String]] = data.map(s => s.trim.split(' '))\n\nval fpg = new FPGrowth()\n.setMinSupport(0.2)\n.setNumPartitions(10)\nval model = fpg.run(transactions)\n\nmodel.freqItemsets.collect().foreach { itemset =>\nprintln(itemset.items.mkString(\"[\", \",\", \"]\") + \", \" + itemset.freq)\n}\n\nval minConfidence = 0.8\nmodel.generateAssociationRules(minConfidence).collect().foreach { rule =>\nprintln(\nrule.antecedent.mkString(\"[\", \",\", \"]\")\n+ \" => \" + rule.consequent .mkString(\"[\", \",\", \"]\")\n+ \", \" + rule.confidence)\n}\n\nFind full example code at \"examples/src/main/scala/org/apache/spark/examples/mllib/SimpleFPGrowth.scala\" in the Spark repo.\n\nFPGrowth implements the FP-growth algorithm. It take an JavaRDD of transactions, where each transaction is an Iterable of items of a generic type. Calling FPGrowth.run with transactions returns an FPGrowthModel that stores the frequent itemsets with their frequencies. The following example illustrates how to mine frequent itemsets and association rules (see Association Rules for details) from transactions.\n\nRefer to the FPGrowth Java docs for details on the API.\n\nimport java.util.Arrays;\nimport java.util.List;\n\nimport org.apache.spark.api.java.JavaRDD;\nimport org.apache.spark.api.java.JavaSparkContext;\n\nimport org.apache.spark.mllib.fpm.AssociationRules;\nimport org.apache.spark.mllib.fpm.FPGrowth;\nimport org.apache.spark.mllib.fpm.FPGrowthModel;\n\nJavaRDD<String> data = sc.textFile(\"data/mllib/sample_fpgrowth.txt\");\n\nJavaRDD<List<String>> transactions = data.map(\nnew Function<String, List<String>>() {\npublic List<String> call(String line) {\nString[] parts = line.split(\" \");\nreturn Arrays.asList(parts);\n}\n}\n);\n\nFPGrowth fpg = new FPGrowth()\n.setMinSupport(0.2)\n.setNumPartitions(10);\nFPGrowthModel<String> model = fpg.run(transactions);\n\nfor (FPGrowth.FreqItemset<String> itemset: model.freqItemsets().toJavaRDD().collect()) {\nSystem.out.println(\"[\" + itemset.javaItems() + \"], \" + itemset.freq());\n}\n\ndouble minConfidence = 0.8;\nfor (AssociationRules.Rule<String> rule\n: model.generateAssociationRules(minConfidence).toJavaRDD().collect()) {\nSystem.out.println(\nrule.javaAntecedent() + \" => \" + rule.javaConsequent() + \", \" + rule.confidence());\n}\n\nFind full example code at \"examples/src/main/java/org/apache/spark/examples/mllib/JavaSimpleFPGrowth.java\" in the Spark repo.\n\nFPGrowth implements the FP-growth algorithm. It take an RDD of transactions, where each transaction is an List of items of a generic type. Calling FPGrowth.train with transactions returns an FPGrowthModel that stores the frequent itemsets with their frequencies.\n\nRefer to the FPGrowth Python docs for more details on the API.\n\nfrom pyspark.mllib.fpm import FPGrowth\n\ndata = sc.textFile(\"data/mllib/sample_fpgrowth.txt\")\ntransactions = data.map(lambda line: line.strip().split(' '))\nmodel = FPGrowth.train(transactions, minSupport=0.2, numPartitions=10)\nresult = model.freqItemsets().collect()\nfor fi in result:\nprint(fi)\n\nFind full example code at \"examples/src/main/python/mllib/fpgrowth_example.py\" in the Spark repo.\n\n## Association Rules\n\nAssociationRules implements a parallel rule generation algorithm for constructing rules that have a single item as the consequent.\n\nRefer to the AssociationRules Scala docs for details on the API.\n\nimport org.apache.spark.mllib.fpm.AssociationRules\nimport org.apache.spark.mllib.fpm.FPGrowth.FreqItemset\n\nval freqItemsets = sc.parallelize(Seq(\nnew FreqItemset(Array(\"a\"), 15L),\nnew FreqItemset(Array(\"b\"), 35L),\nnew FreqItemset(Array(\"a\", \"b\"), 12L)\n))\n\nval ar = new AssociationRules()\n.setMinConfidence(0.8)\nval results = ar.run(freqItemsets)\n\nresults.collect().foreach { rule =>\nprintln(\"[\" + rule.antecedent.mkString(\",\")\n+ \"=>\"\n+ rule.consequent.mkString(\",\") + \"],\" + rule.confidence)\n}\n\nFind full example code at \"examples/src/main/scala/org/apache/spark/examples/mllib/AssociationRulesExample.scala\" in the Spark repo.\n\nAssociationRules implements a parallel rule generation algorithm for constructing rules that have a single item as the consequent.\n\nRefer to the AssociationRules Java docs for details on the API.\n\nimport java.util.Arrays;\n\nimport org.apache.spark.api.java.JavaRDD;\nimport org.apache.spark.api.java.JavaSparkContext;\nimport org.apache.spark.mllib.fpm.AssociationRules;\nimport org.apache.spark.mllib.fpm.FPGrowth;\nimport org.apache.spark.mllib.fpm.FPGrowth.FreqItemset;\n\nJavaRDD<FPGrowth.FreqItemset<String>> freqItemsets = sc.parallelize(Arrays.asList(\nnew FreqItemset<String>(new String[] {\"a\"}, 15L),\nnew FreqItemset<String>(new String[] {\"b\"}, 35L),\nnew FreqItemset<String>(new String[] {\"a\", \"b\"}, 12L)\n));\n\nAssociationRules arules = new AssociationRules()\n.setMinConfidence(0.8);\nJavaRDD<AssociationRules.Rule<String>> results = arules.run(freqItemsets);\n\nfor (AssociationRules.Rule<String> rule : results.collect()) {\nSystem.out.println(\nrule.javaAntecedent() + \" => \" + rule.javaConsequent() + \", \" + rule.confidence());\n}\n\nFind full example code at \"examples/src/main/java/org/apache/spark/examples/mllib/JavaAssociationRulesExample.java\" in the Spark repo.\n\n## PrefixSpan\n\nPrefixSpan is a sequential pattern mining algorithm described in Pei et al., Mining Sequential Patterns by Pattern-Growth: The PrefixSpan Approach. We refer the reader to the referenced paper for formalizing the sequential pattern mining problem.\n\nspark.mllib’s PrefixSpan implementation takes the following parameters:\n\n• minSupport: the minimum support required to be considered a frequent sequential pattern.\n• maxPatternLength: the maximum length of a frequent sequential pattern. Any frequent pattern exceeding this length will not be included in the results.\n• maxLocalProjDBSize: the maximum number of items allowed in a prefix-projected database before local iterative processing of the projected database begins. This parameter should be tuned with respect to the size of your executors.\n\nExamples\n\nThe following example illustrates PrefixSpan running on the sequences (using same notation as Pei et al):\n\n <(12)3>\n<1(32)(12)>\n<(12)5>\n<6>\n\n\nPrefixSpan implements the PrefixSpan algorithm. Calling PrefixSpan.run returns a PrefixSpanModel that stores the frequent sequences with their frequencies.\n\nRefer to the PrefixSpan Scala docs and PrefixSpanModel Scala docs for details on the API.\n\nimport org.apache.spark.mllib.fpm.PrefixSpan\n\nval sequences = sc.parallelize(Seq(\nArray(Array(1, 2), Array(3)),\nArray(Array(1), Array(3, 2), Array(1, 2)),\nArray(Array(1, 2), Array(5)),\nArray(Array(6))\n), 2).cache()\nval prefixSpan = new PrefixSpan()\n.setMinSupport(0.5)\n.setMaxPatternLength(5)\nval model = prefixSpan.run(sequences)\nmodel.freqSequences.collect().foreach { freqSequence =>\nprintln(\nfreqSequence.sequence.map(_.mkString(\"[\", \", \", \"]\")).mkString(\"[\", \", \", \"]\") +\n\", \" + freqSequence.freq)\n}\n\nFind full example code at \"examples/src/main/scala/org/apache/spark/examples/mllib/PrefixSpanExample.scala\" in the Spark repo.\n\nPrefixSpan implements the PrefixSpan algorithm. Calling PrefixSpan.run returns a PrefixSpanModel that stores the frequent sequences with their frequencies.\n\nRefer to the PrefixSpan Java docs and PrefixSpanModel Java docs for details on the API.\n\nimport java.util.Arrays;\nimport java.util.List;\n\nimport org.apache.spark.mllib.fpm.PrefixSpan;\nimport org.apache.spark.mllib.fpm.PrefixSpanModel;\n\nJavaRDD<List<List<Integer>>> sequences = sc.parallelize(Arrays.asList(\nArrays.asList(Arrays.asList(1, 2), Arrays.asList(3)),\nArrays.asList(Arrays.asList(1), Arrays.asList(3, 2), Arrays.asList(1, 2)),\nArrays.asList(Arrays.asList(1, 2), Arrays.asList(5)),\nArrays.asList(Arrays.asList(6))\n), 2);\nPrefixSpan prefixSpan = new PrefixSpan()\n.setMinSupport(0.5)\n.setMaxPatternLength(5);\nPrefixSpanModel<Integer> model = prefixSpan.run(sequences);\nfor (PrefixSpan.FreqSequence<Integer> freqSeq: model.freqSequences().toJavaRDD().collect()) {\nSystem.out.println(freqSeq.javaSequence() + \", \" + freqSeq.freq());\n}\n\nFind full example code at \"examples/src/main/java/org/apache/spark/examples/mllib/JavaPrefixSpanExample.java\" in the Spark repo." ]
[ null ]
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https://segmentfault.com/a/1190000041330561
[ "# 每日一练(11):调整数组顺序使奇数位于偶数前面\n\ntitle: 每日一练(11):调整数组顺序使奇数位于偶数前面\n\ncategories:[剑指offer]\n\ntags:[每日一练]\n\ndate: 2022/01/24\n\n## 每日一练(11):调整数组顺序使奇数位于偶数前面\n\n0 <= nums.length <= 50000\n\n0 <= nums[i] <= 10000\n\n### 方法一:有手就行\n\n``````vector<int> exchange(vector<int>& nums) {\nvector<int>res;\nfor (int i = 0; i < nums.size(); i++) {\nif ((nums[i] & 1) != 0) { //nums[i] & 1) != 0 可用 nums[i] % 2 替代\nres.push_back(nums[i]); //放入奇数\n}\n}\nfor (int i = 0; i < nums.size(); i++) {\nif ((nums[i] & 1) != 1) { //nums[i] & 1) != 1 可用 (nums[i] % 2) == 0 替代\nres.push_back(nums[i]); //放入偶数\n}\n}\nreturn res;\n}``````\n\n### 方法二:首尾双指针\n\n• 定义头指针 left ,尾指针 right .\n• left 一直往右移,直到它指向的值为偶数 .\n• right 一直往左移, 直到它指向的值为奇数 .\n• 交换 nums[left] 和 nums[right] .\n• 重复上述操作,直到 left == right .\n``````vector<int> exchange(vector<int>& nums) {\nint left = 0, right = nums.size() - 1;\nwhile (left < right) {\nif ((nums[left] & 1) != 0) { //如果为偶数\nleft++;\ncontinue;\n}\nif ((nums[right] & 1) != 1) { //如果为奇数\nright--;\ncontinue;\n}\nswap(nums[left++], nums[right--]); //交换值\n}\nreturn nums;\n}``````\n\n### 方法三:快慢双指针\n\n• 定义快慢双指针 fast 和 slow ,fast 在前, slow 在后 .\n• fast 的作用是向前搜索奇数位置,slow 的作用是指向下一个奇数应当存放的位置 .\n• fast 向前移动,当它搜索到奇数时,将它和 nums[slow ] 交换,此时 slow 向前移动一个位置 .\n• 重复上述操作,直到 fast 指向数组末尾 .\n``````vector<int> exchange(vector<int>& nums) {\nint slow = 0, fast = 0;\nwhile (fast < nums.size()) {\nif ((nums[fast]) & 1) { //奇数\nswap(nums[slow], nums[fast]);//交换\nslow++;\n}\nfast++;\n}\nreturn nums;\n}``````\n\n25 声望\n7 粉丝\n0 条评论", null, "" ]
[ null, "https://image-static.segmentfault.com/317/931/3179314346-5f61e47221e07", null ]
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http://www.stock-charts-made-easy.com/fibonacci.html
[ "## \"Can Fibonacci Price Projections Forecast Key Turning Points In The Markets?...\"", null, "es... and, with great accuracy. You'll see how calculations based on the Fibonacci number series can be applied to stock prices for startling results.", null, "", null, "", null, "", null, "Applied correctly, support and resistance levels will materialize. Likely turning points will present themselves, clearly.\n\nSo, how are the retracement percentages derived?\n\nBegin by adding one to itself... which is two. Then add two to one to get three -- the next number in the series. You keep adding the current number to the previous number so it looks like this...\n\n1, 1, 2, 3, 5, 8, 13, 21, 34, 55, 89, 144...\n\nNow, you have a number series which goes on forever.", null, "Next, there is an important ratio which is calculated by dividing one number in the series by its previous number. For example...\n\n144 ÷ 89 = 1.618\n\nThis ratio is known as the Golden Mean and is the basis for most retracement percentages. Equally important is the reciprocal of 1.618...\n\n1 ÷ 1.618 = .618\n\nHere are some other percentages you'll want to use...\n\n.382 = .618 squared\n\n.500 = 1 ÷ 2, the second and third numbers in the series\n\n.786 = square root of .618\n\n1.000 = 1.618 x .618\n\n1.272 = square root of 1.618\n\n2.618 = 1.618 squared", null, "### \"How are Fibonacci percentages applied?\"\n\nBegin by locating isolated highs and lows in a trend. The various swing points will be used to measure the retracements.\n\nIn the following illustration, trend AB has been established...", null, "Counter-trend BC is a 50-percent retracement of trend AB. Trend CD continues in the same direction as trend AB. Counter-trend DE is a 38-percent retracement of trend CD.\n\nShallow retracements are typical in a strong trend.", null, "Combining retracements from multiple swing points is a dynamic forecasting method. In the next example, retracement percentages are calculated from trends AD and CD...", null, "The 38-percent retracement is calculated from trend AD. The 50-percent retracement is from trend CD. When retracements from different swing points cluster together it is known as confluence.\n\nAreas of confluence are likely turning points.\n\nClick For Fibonacci Confluence Chart Examples", null, "So far, you have seen how to calculate probable retracements using the high and low swing point of an established trend. You can use Fibonacci percentages, in combination, to forecast the end of a trend... with excellent results.\n\nThis next example shows a typical five-wave trend progression where the 1.618 and .618 percentage extensions are used...", null, "First, trend AB is multiplied by 1.618 and then added to itself... forecasting a target price at point F. Next, counter-trend BC retraces part of trend AB before continuing on. Finally, trend AD is multiplied by .618 and added to itself. This projection also forecasts the trend termination at point F. Incredibly powerful!\n\nIn a five-wave pattern, this ratio combination will bear-out turning points, again and again.", null, "There are many other ways to apply Fibonacci retracements. Purchase the following books to learn more...", null, "by Robert Fischer\nJohn Wiley & Sons, Inc.\n2001\n352 Pages\n\nGet This Book, Now!", null, "Dynamic Trading - Dynamic Concepts in Time, Price and Pattern Analysis with Practical Strategies for Traders and Investors\nby Robert C. Miner\n1997\n564 Pages\n\n Or... SEARCH for a BOOK Enter title, author, item# or ISBN", null, "Home | Charting Principles | Forecasting Tools | Trading Systems | Market Metrics\n\nCharts In Action | Market Commentary | Online Newsletter | Bookstore\n\nFor the Record | Site Map", null, "", null, "", null, "Online Newsletter Subscribe FREE! And... receive your FREE report. Back Issues ofChart Wealth", null, "What's an RSS feed?", null, "", null, "", null, "", null, "", null, "" ]
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https://se.mathworks.com/help/simulink/slref/discretepidcontroller2dof.html
[ "# Discrete PID Controller (2DOF)\n\nDiscrete-time or continuous-time two-degree-of-freedom PID controller\n\n• Library:\n\n•", null, "## Description\n\nThe Discrete PID Controller (2DOF) block implements a two-degree-of-freedom PID controller (PID, PI, or PD). The block is identical to the PID Controller (2DOF) block with the Time domain parameter set to `Discrete-time`.\n\nThe block generates an output signal based on the difference between a reference signal and a measured system output. The block computes a weighted difference signal for the proportional and derivative actions according to the setpoint weights (b and c) that you specify. The block output is the sum of the proportional, integral, and derivative actions on the respective difference signals, where each action is weighted according to the gain parameters P, I, and D. A first-order pole filters the derivative action.\n\nThe block supports several controller types and structures. Configurable options in the block include:\n\n• Controller type (PID, PI, or PD) — See the Controller parameter.\n\n• Controller form (Parallel or Ideal) — See the Form parameter.\n\n• Time domain (discrete or continuous) — See the Time domain parameter.\n\n• Initial conditions and reset trigger — See the Source and External reset parameters.\n\n• Output saturation limits and built-in anti-windup mechanism — See the Limit output parameter.\n\n• Signal tracking for bumpless control transfer and multiloop control — See the Enable tracking mode parameter.\n\nAs you change these options, the internal structure of the block changes by activating different variant subsystems. (See Variant Subsystems.) To examine the internal structure of the block and its variant subsystems, right-click the block and select Mask > Look Under Mask.\n\n### Control Configuration\n\nIn one common implementation, the PID Controller block operates in the feedforward path of a feedback loop.", null, "For a single-input block that accepts an error signal (a difference between a setpoint and a system output), see Discrete PID Controller.\n\n### PID Gain Tuning\n\nThe PID controller coefficients and the setpoint weights are tunable either manually or automatically. Automatic tuning requires Simulink® Control Design™ software. For more information about automatic tuning, see the Select tuning method parameter.\n\n## Ports\n\n### Input\n\nexpand all\n\nReference signal for plant to follow, as shown.", null, "When the reference signal is a vector, the block acts separately on each signal, vectorizing the PID coefficients and producing a vector output signal of the same dimensions. You can specify the PID coefficients and some other parameters as vectors of the same dimensions as the input signal. Doing so is equivalent to specifying a separate PID controller for each entry in the input signal.\n\nData Types: `single` | `double` | `int8` | `int16` | `int32` | `int64` | `uint8` | `uint16` | `uint32` | `uint64` | `fixed point`\n\nFeedback signal for the controller, from the plant output.\n\nData Types: `single` | `double` | `int8` | `int16` | `int32` | `int64` | `uint8` | `uint16` | `uint32` | `uint64` | `fixed point`\n\nProportional gain, provided from a source external to the block. External gain input is useful, for example, when you want to map a different PID parameterization to the PID gains of the block. You can also use external gain input to implement gain-scheduled PID control. In gain-scheduled control, you determine the PID coefficients by logic or other calculation in your model and feed them to the block.\n\n#### Dependencies\n\nTo enable this port, set Controller parameters Source to `external`.\n\nData Types: `single` | `double` | `int8` | `int16` | `int32` | `int64` | `uint8` | `uint16` | `uint32` | `uint64` | `fixed point`\n\nIntegral gain, provided from a source external to the block. External gain input is useful, for example, when you want to map a different PID parameterization to the PID gains of the block. You can also use external gain input to implement gain-scheduled PID control. In gain-scheduled control, you determine the PID coefficients by logic or other calculation in your model and feed them to the block.\n\nWhen you supply gains externally, time variations in the integral gain are also integrated. This result occurs because of the way the PID gains are implemented within the block. For details, see the Controller parameters Source parameter.\n\n#### Dependencies\n\nTo enable this port, set Controller parameters Source to `external`, and set Controller to a controller type that has integral action.\n\nData Types: `single` | `double` | `int8` | `int16` | `int32` | `int64` | `uint8` | `uint16` | `uint32` | `uint64` | `fixed point`\n\nDerivative gain, provided from a source external to the block. External gain input is useful, for example, when you want to map a different PID parameterization to the PID gains of the block. You can also use external gain input to implement gain-scheduled PID control. In gain-scheduled control, you determine the PID coefficients by logic or other calculation in your model and feed them to the block.\n\nWhen you supply gains externally, time variations in the derivative gain are also differentiated. This result occurs because of the way the PID gains are implemented within the block. For details, see the Controller parameters Source parameter.\n\n#### Dependencies\n\nTo enable this port, set Controller parameters Source to `external`, and set Controller to a controller type that has derivative action.\n\nData Types: `single` | `double` | `int8` | `int16` | `int32` | `int64` | `uint8` | `uint16` | `uint32` | `uint64` | `fixed point`\n\nDerivative filter coefficient, provided from a source external to the block. External coefficient input is useful, for example, when you want to map a different PID parameterization to the PID gains of the block. You can also use the external input to implement gain-scheduled PID control. In gain-scheduled control, you determine the PID coefficients by logic or other calculation in your model and feed them to the block.\n\n#### Dependencies\n\nTo enable this port, set Controller parameters Source to `external`, and set Controller to a controller type that has a filtered derivative.\n\nData Types: `single` | `double` | `int8` | `int16` | `int32` | `int64` | `uint8` | `uint16` | `uint32` | `uint64` | `fixed point`\n\nProportional setpoint weight, provided from a source external to the block. External input is useful, for example, when you want to map a different PID parameterization to the PID gains of the block. You can also use the external input to implement gain-scheduled PID control. In gain-scheduled control, you determine the PID coefficients by logic or other calculation in your model and feed them to the block.\n\n#### Dependencies\n\nTo enable this port, set Controller parameters Source to `external`.\n\nData Types: `single` | `double` | `int8` | `int16` | `int32` | `int64` | `uint8` | `uint16` | `uint32` | `uint64` | `fixed point`\n\nDerivative setpoint weight, provided from a source external to the block. External input is useful, for example, when you want to map a different PID parameterization to the PID gains of the block. You can also use the external input to implement gain-scheduled PID control. In gain-scheduled control, you determine the PID coefficients by logic or other calculation in your model and feed them to the block.\n\n#### Dependencies\n\nTo enable this port, set Controller parameters Source to `external`, and set Controller to a controller type that has derivative action.\n\nData Types: `single` | `double` | `int8` | `int16` | `int32` | `int64` | `uint8` | `uint16` | `uint32` | `uint64` | `fixed point`\n\nTrigger to reset the integrator and filter to their initial conditions. Use the External reset parameter to specify what kind of signal triggers a reset. The port icon indicates the trigger type specified in that parameter. For example, the following illustration shows a continuous-time PID Controller (2DOF) block with External reset set to `rising`.", null, "When the trigger occurs, the block resets the integrator and filter to the initial conditions specified by the Integrator Initial condition and Filter Initial condition parameters or the I0 and D0 ports.\n\n### Note\n\nTo be compliant with the Motor Industry Software Reliability Association (MISRA®) software standard, your model must use Boolean signals to drive the external reset ports of the PID controller block.\n\n#### Dependencies\n\nTo enable this port, set External reset to any value other than `none`.\n\nData Types: `single` | `double` | `int8` | `int16` | `int32` | `int64` | `uint8` | `uint16` | `uint32` | `uint64` | `fixed point` | `Boolean`\n\nIntegrator initial condition, provided from a source external to the block.\n\n#### Dependencies\n\nTo enable this port, set Initial conditions Source to `external`, and set Controller to a controller type that has integral action.\n\nData Types: `single` | `double` | `int8` | `int16` | `int32` | `int64` | `uint8` | `uint16` | `uint32` | `uint64` | `fixed point`\n\nInitial condition of the derivative filter, provided from a source external to the block.\n\n#### Dependencies\n\nTo enable this port, set Initial conditions Source to `external`, and set Controller to a controller type that has derivative action.\n\nData Types: `single` | `double` | `int8` | `int16` | `int32` | `int64` | `uint8` | `uint16` | `uint32` | `uint64` | `fixed point`\n\nSignal for controller output to track. When signal tracking is active, the difference between the tracking signal and the block output is fed back to the integrator input. Signal tracking is useful for implementing bumpless control transfer in systems that switch between two controllers. It can also be useful to prevent block windup in multiloop control systems. For more information, see the Enable tracking mode parameter.\n\n#### Dependencies\n\nTo enable this port, select the Enable tracking mode parameter.\n\nData Types: `single` | `double` | `int8` | `int16` | `int32` | `int64` | `uint8` | `uint16` | `uint32` | `uint64` | `fixed point`\n\nDiscrete-integrator time, provided as a scalar to the block. You can use your own value of discrete-time integrator sample time that defines the rate at which the block is going to be run either in Simulink or on external hardware. The value of the discrete-time integrator time should match the average sampling rate of the external interrupts, when the block is used inside a conditionally-executed subsystem.\n\nIn other words, you can specify `Ts` for any of the integrator methods below such that the value matches the average sampling rate of the external interrupts. In discrete time, the derivative term of the controller transfer function is:\n\n`$D\\left[\\frac{N}{1+N\\alpha \\left(z\\right)}\\right],$`\n\nwhere α(z) depends on the integrator method you specify with this parameter.\n\n```Forward Euler```\n`$\\alpha \\left(z\\right)=\\frac{{T}_{s}}{z-1}.$`\n```Backward Euler```\n`$\\alpha \\left(z\\right)=\\frac{{T}_{s}z}{z-1}.$`\n`Trapezoidal`\n`$\\alpha \\left(z\\right)=\\frac{{T}_{s}}{2}\\frac{z+1}{z-1}.$`\n\n#### Dependencies\n\nTo enable this port, set Time Domain to `Discrete-time` and select the PID Controller is inside a conditionally executed subsystem option.\n\nData Types: `single` | `double` | `int8` | `int16` | `int32` | `int64` | `uint8` | `uint16` | `uint32` | `uint64`\n\n### Output\n\nexpand all\n\nController output, generally based on a sum of the input signal, the integral of the input signal, and the derivative of the input signal, weighted by the setpoint weights and by the proportional, integral, and derivative gain parameters. A first-order pole filters the derivative action. Which terms are present in the controller signal depends on what you select for the Controller parameter. The base controller transfer function for the current settings is displayed in the Compensator formula section of the block parameters and under the mask. Other parameters modify the block output, such as saturation limits specified by the Upper Limit and Lower Limit saturation parameters.\n\nThe controller output is a vector signal when any of the inputs is a vector signal. In that case, the block acts as N independent PID controllers, where N is the number of signals in the input vector.\n\nData Types: `single` | `double` | `int8` | `int16` | `int32` | `int64` | `uint8` | `uint16` | `uint32` | `uint64` | `fixed point`\n\n## Parameters\n\nexpand all\n\nSpecify which of the proportional, integral, and derivative terms are in the controller.\n\n`PID `\n\nProportional, integral, and derivative action.\n\n`PI`\n\nProportional and integral action only.\n\n`PD`\n\nProportional and derivative action only.\n\n### Tip\n\nThe controller output for the current setting is displayed in the Compensator formula section of the block parameters and under the mask.\n\n#### Programmatic Use\n\n Block Parameter: `Controller` Type: string, character vector Values: `\"PID\"`, `\"PI\"`, `\"PD\"` Default: `\"PID\"`\n\nSpecify whether the controller structure is parallel or ideal.\n\n`Parallel`\n\nThe proportional, integral, and derivative gains P, I, and D, are applied independently. For example, for a continuous-time 2-DOF PID controller in parallel form, the controller output u is:\n\n`$u=P\\left(br-y\\right)+I\\frac{1}{s}\\left(r-y\\right)+D\\frac{N}{1+N\\frac{1}{s}}\\left(cr-y\\right),$`\n\nwhere r is the reference signal, y is the measured plant output signal, and b and c are the setpoint weights.\n\nFor a discrete-time 2-DOF controller in parallel form, the controller output is:\n\n`$u=P\\left(br-y\\right)+I\\alpha \\left(z\\right)\\left(r-y\\right)+D\\frac{N}{1+N\\beta \\left(z\\right)}\\left(cr-y\\right),$`\n\nwhere the Integrator method and Filter method parameters determine α(z) and β(z), respectively.\n\n`Ideal`\n\nThe proportional gain P acts on the sum of all actions. For example, for a continuous-time 2-DOF PID controller in ideal form, the controller output is:\n\n`$u=P\\left[\\left(br-y\\right)+I\\frac{1}{s}\\left(r-y\\right)+D\\frac{N}{1+N\\frac{1}{s}}\\left(cr-y\\right)\\right].$`\n\nFor a discrete-time 2-DOF PID controller in ideal form, the transfer function is:\n\n`$u=P\\left[\\left(br-y\\right)+I\\alpha \\left(z\\right)\\left(r-y\\right)+D\\frac{N}{1+N\\beta \\left(z\\right)}\\left(cr-y\\right)\\right],$`\n\nwhere the Integrator method and Filter method parameters determine α(z) and β(z), respectively.\n\n### Tip\n\nThe controller output for the current settings is displayed in the Compensator formula section of the block parameters and under the mask.\n\n#### Programmatic Use\n\n Block Parameter: `Controller` Type: string, character vector Values: `\"Parallel\"`, `\"Ideal\"` Default: `\"Parallel\"`\n\nWhen you select `Discrete-time`, it is recommended that you specify an explicit sample time for the block. See the Sample time (-1 for inherited) parameter. Selecting `Discrete-time` also enables the Integrator method, and Filter method parameters.\n\nWhen the PID Controller block is in a model with synchronous state control (see the State Control block), you cannot select `Continuous-time`.\n\n### Note\n\nThe PID Controller (2DOF) and Discrete PID Controller (2DOF) blocks are identical except for the default value of this parameter.\n\n#### Programmatic Use\n\n Block Parameter: `TimeDomain` Type: string, character vector Values: `\"Continuous-time\"`, `\"Discrete-time\"` Default: `\"Discrete-time\"`\n\nFor discrete-time PID controllers, enable the discrete-time integrator port to use your own value of discrete-time integrator sample time. To ensure proper integration, Use the `TDTI` port to provide a scalar value of Δt for accurate discrete-time integration.\n\n#### Dependencies\n\nTo enable this parameter, set Time Domain to `Discrete-time`.\n\n#### Programmatic Use\n\n Block Parameter: `UseExternalTs` Type: string, character vector Values: `\"on\"`, `\"off\"` Default: `\"off\"`\n\nSpecify a sample time by entering a positive scalar value, such as 0.1. The default discrete sample time of –1 means that the block inherits its sample time from upstream blocks. However, it is recommended that you set the controller sample time explicitly, especially if you expect the sample time of upstream blocks to change. The effect of the controller coefficients P, I, D, and N depend on the sample time. Thus, for a given set of coefficient values, changing the sample time changes the performance of the controller.\n\nTo implement a continuous-time controller, set Time domain to `Continuous-time`.\n\n### Tip\n\nIf you want to run the block with an externally specified or variable sample time, set this parameter to –1 and put the block in a Triggered Subsystem. Then, trigger the subsystem at the desired sample time.\n\n#### Dependencies\n\nTo enable this parameter, set Time domain to `Discrete-time`.\n\n#### Programmatic Use\n\n Block Parameter: `SampleTime` Type: scalar Values: `-1`, positive scalar Default: `-1`\n\nIn discrete time, the integral term of the controller transfer function is Ia(z), where a(z) depends on the integrator method you specify with this parameter.\n\n`Forward Euler`\n\nForward rectangular (left-hand) approximation,\n\n`$\\alpha \\left(z\\right)=\\frac{{T}_{s}}{z-1}.$`\n\nThis method is best for small sampling times, where the Nyquist limit is large compared to the bandwidth of the controller. For larger sampling times, the `Forward Euler` method can result in instability, even when discretizing a system that is stable in continuous time.\n\n`Backward Euler`\n\nBackward rectangular (right-hand) approximation,\n\n`$\\alpha \\left(z\\right)=\\frac{{T}_{s}z}{z-1}.$`\n\nAn advantage of the `Backward Euler` method is that discretizing a stable continuous-time system using this method always yields a stable discrete-time result.\n\n`Trapezoidal`\n\nBilinear approximation,\n\n`$\\alpha \\left(z\\right)=\\frac{{T}_{s}}{2}\\frac{z+1}{z-1}.$`\n\nAn advantage of the `Trapezoidal` method is that discretizing a stable continuous-time system using this method always yields a stable discrete-time result. Of all available integration methods, the `Trapezoidal` method yields the closest match between frequency-domain properties of the discretized system and the corresponding continuous-time system.\n\n### Tip\n\nThe controller formula for the current setting is displayed in the Compensator formula section of the block parameters and under the mask.\n\n#### Dependencies\n\nTo enable this parameter, set Time Domain to `Discrete-time` and set Controller to a controller type with integral action.\n\n#### Programmatic Use\n\n Block Parameter: `IntegratorMethod` Type: string, character vector Values: `\"Forward Euler\"`, `\"Backward Euler\"`, `\"Trapezoidal\"` Default: `\"Forward Euler\"`\n\nIn discrete time, the derivative term of the controller transfer function is:\n\n`$D\\left[\\frac{N}{1+N\\alpha \\left(z\\right)}\\right],$`\n\nwhere α(z) depends on the filter method you specify with this parameter.\n\n`Forward Euler`\n\nForward rectangular (left-hand) approximation,\n\n`$\\alpha \\left(z\\right)=\\frac{{T}_{s}}{z-1}.$`\n\nThis method is best for small sampling times, where the Nyquist limit is large compared to the bandwidth of the controller. For larger sampling times, the `Forward Euler` method can result in instability, even when discretizing a system that is stable in continuous time.\n\n`Backward Euler`\n\nBackward rectangular (right-hand) approximation,\n\n`$\\alpha \\left(z\\right)=\\frac{{T}_{s}z}{z-1}.$`\n\nAn advantage of the `Backward Euler` method is that discretizing a stable continuous-time system using this method always yields a stable discrete-time result.\n\n`Trapezoidal`\n\nBilinear approximation,\n\n`$\\alpha \\left(z\\right)=\\frac{{T}_{s}}{2}\\frac{z+1}{z-1}.$`\n\nAn advantage of the `Trapezoidal` method is that discretizing a stable continuous-time system using this method always yields a stable discrete-time result. Of all available integration methods, the `Trapezoidal` method yields the closest match between frequency-domain properties of the discretized system and the corresponding continuous-time system.\n\n### Tip\n\nThe controller formula for the current setting is displayed in the Compensator formula section of the block parameters and under the mask.\n\n#### Dependencies\n\nTo enable this parameter, set Time Domain to `Discrete-time` and enable Use filtered derivative.\n\n#### Programmatic Use\n\n Block Parameter: `FilterMethod` Type: string, character vector Values: `\"Forward Euler\"`, `\"Backward Euler\"`, `\"Trapezoidal\"` Default: `\"Forward Euler\"`\n\n### Main\n\n`internal`\n\nSpecify the controller gains, filter coefficient, and setpoint weights using the block parameters P, I, D, N, b, and c respectively.\n\n`external`\n\nSpecify the PID gains, filter coefficient, and setpoint weights externally using block inputs. An additional input port appears on the block for each parameter that is required for the current controller type.\n\nEnabling external inputs for the parameters allows you to compute their values externally to the block and provide them to the block as signal inputs.\n\nExternal input is useful, for example, when you want to map a different PID parameterization to the PID gains of the block. You can also use external gain input to implement gain-scheduled PID control. In gain-scheduled control, you determine the PID gains by logic or other calculation in your model and feed them to the block.\n\nWhen you supply gains externally, time variations in the integral and derivative gain values are integrated and differentiated, respectively. The derivative setpoint weight c is also differentiated. This result occurs because in both continuous time and discrete time, the gains are applied to the signal before integration or differentiation. For example, for a continuous-time PID controller with external inputs, the integrator term is implemented as shown in the following illustration.", null, "Within the block, the input signal u is multiplied by the externally supplied integrator gain, I, before integration. This implementation yields:\n\n`${u}_{i}=\\int \\left(r-y\\right)I\\text{\\hspace{0.17em}}dt.$`\n\nThus, the integrator gain is included in the integral. Similarly, in the derivative term of the block, multiplication by the derivative gain precedes the differentiation, which causes the derivative gain D and the derivative setpoint weight c to be differentiated.\n\n#### Programmatic Use\n\n Block Parameter: `ControllerParametersSource` Type: string, character vector Values: `\"internal\"`, `\"external\"` Default: `\"internal\"`\n\nSpecify a finite, real gain value for the proportional gain. When Controller form is:\n\n• `Parallel` — Proportional action is independent of the integral and derivative actions. For example, for a continuous-time 2-DOF PID controller in parallel form, the controller output u is:\n\n`$u=P\\left(br-y\\right)+I\\frac{1}{s}\\left(r-y\\right)+D\\frac{N}{1+N\\frac{1}{s}}\\left(cr-y\\right),$`\n\nwhere r is the reference signal, y is the measured plant output signal, and b and c are the setpoint weights.\n\nFor a discrete-time 2-DOF controller in parallel form, the controller output is:\n\n`$u=P\\left(br-y\\right)+I\\alpha \\left(z\\right)\\left(r-y\\right)+D\\frac{N}{1+N\\beta \\left(z\\right)}\\left(cr-y\\right),$`\n\nwhere the Integrator method and Filter method parameters determine α(z) and β(z), respectively.\n\n• `Ideal` — The proportional gain multiples the integral and derivative terms. For example, for a continuous-time 2-DOF PID controller in ideal form, the controller output is:\n\n`$u=P\\left[\\left(br-y\\right)+I\\frac{1}{s}\\left(r-y\\right)+D\\frac{N}{1+N\\frac{1}{s}}\\left(cr-y\\right)\\right].$`\n\nFor a discrete-time 2-DOF PID controller in ideal form, the transfer function is:\n\n`$u=P\\left[\\left(br-y\\right)+I\\alpha \\left(z\\right)\\left(r-y\\right)+D\\frac{N}{1+N\\beta \\left(z\\right)}\\left(cr-y\\right)\\right],$`\n\nwhere the Integrator method and Filter method parameters determine α(z) and β(z), respectively.\n\nTunable: Yes\n\n#### Dependencies\n\nTo enable this parameter, set the Controller parameters Source to `internal`.\n\n#### Programmatic Use\n\n Block Parameter: `P` Type: scalar, vector Default: 1\n\nSpecify a finite, real gain value for the integral gain.\n\nTunable: Yes\n\n#### Dependencies\n\nTo enable this parameter, in the Main tab, set the controller-parameters Source to `internal`, and set Controller to a type that has integral action.\n\n#### Programmatic Use\n\n Block Parameter: `I` Type: scalar, vector Default: 1\n\nSpecify a finite, real gain value for the derivative gain.\n\nTunable: Yes\n\n#### Dependencies\n\nTo enable this parameter, in the Main tab, set the controller-parameters Source to `internal`, and set Controller to `PID` or `PD`.\n\n#### Programmatic Use\n\n Block Parameter: `D` Type: scalar, vector Default: 0\n\nFor discrete-time PID controllers only, clear this option to replace the filtered derivative with an unfiltered discrete-time differentiator. When you do so, the derivative term of the controller output becomes:\n\n`$D\\frac{z-1}{z{T}_{s}}\\left(cr-y\\right).$`\n\nFor continuous-time PID controllers, the derivative term is always filtered.\n\n#### Dependencies\n\nTo enable this parameter, set Time domain to `Discrete-time`, and set Controller to a type that has a derivative term.\n\n#### Programmatic Use\n\n Block Parameter: `UseFilter` Type: string, character vector Values: `\"on\"`, `\"off\"` Default: `\"on\"`\n\nSpecify a finite, real gain value for the filter coefficient. The filter coefficient determines the pole location of the filter in the derivative action of the block. The location of the filter pole depends on the Time domain parameter.\n\n• When Time domain is `Continuous-time`, the pole location is `s = -N`.\n\n• When Time domain is `Discrete-time`, the pole location depends on the Filter method parameter.\n\nFilter MethodLocation of Filter Pole\n```Forward Euler```${z}_{pole}=1-N{T}_{s}$\n```Backward Euler```${z}_{pole}=\\frac{1}{1+N{T}_{s}}$\n`Trapezoidal`${z}_{pole}=\\frac{1-N{T}_{s}/2}{1+N{T}_{s}/2}$\n\nThe block does not support `N = Inf` (ideal unfiltered derivative). When the Time domain is `Discrete-time`, you can clear Use filtered derivative to remove the derivative filter.\n\nTunable: Yes\n\n#### Dependencies\n\nTo enable this parameter, in the Main tab, set the controller-parameters Source to `internal` and set Controller to `PID` or `PD`.\n\n#### Programmatic Use\n\n Block Parameter: `N` Type: scalar, vector Default: 100\n\nSetpoint weight on the proportional term of the controller. The proportional term of a 2-DOF controller output is P(bry), where r is the reference signal and y is the measured plant output. Setting b to 0 eliminates proportional action on the reference signal, which can reduce overshoot in the system response to step changes in the setpoint. Changing the relative values of b and c changes the balance between disturbance rejection and setpoint tracking.\n\nTunable: Yes\n\n#### Dependencies\n\nTo enable this parameter, in the Main tab, set the controller-parameters Source to `internal`.\n\n#### Programmatic Use\n\n Block Parameter: `b` Type: scalar, vector Default: 1\n\nSetpoint weight on the derivative term of the controller. The derivative term of a 2-DOF controller acts on cry, where r is the reference signal and y is the measured plant output. Thus, setting c to 0 eliminates derivative action on the reference signal, which can reduce transient response to step changes in the setpoint. Setting c to 0 can yield a controller that achieves both effective disturbance rejection and smooth setpoint tracking without excessive transient response. Changing the relative values of b and c changes the balance between disturbance rejection and setpoint tracking.\n\nTunable: Yes\n\n#### Dependencies\n\nTo enable this parameter, in the Main tab, set the controller-parameters Source to `internal` and set Controller to a type that has derivative action.\n\n#### Programmatic Use\n\n Block Parameter: `c` Type: scalar, vector Default: 1\n\nIf you have Simulink Control Design software, you can automatically tune the PID coefficients when they are internal to the block. To do so, use this parameter to select a tuning tool, and click Tune.\n\n`Transfer Function Based (PID Tuner App)`\n\nUse PID Tuner, which lets you interactively tune PID coefficients while examining relevant system responses to validate performance. PID Tuner can tune all the coefficients P, I, D, and N, and the setpoint coefficients b and c. By default, PID Tuner works with a linearization of your plant model. For models that cannot be linearized, you can tune PID coefficients against a plant model estimated from simulated or measured response data. For more information, see Design Two-Degree-of-Freedom PID Controllers (Simulink Control Design).\n\n`Frequency Response Based`\n\nUse Frequency Response Based PID Tuner, which tunes PID controller coefficients based on frequency-response estimation data obtained by simulation. This tuning approach is especially useful for plants that are not linearizable or that linearize to zero. Frequency Response Based PID Tuner tunes the coefficients P, I, D, and N, but does not tune the setpoint coefficients b and c. For more information, see Design PID Controller from Plant Frequency-Response Data (Simulink Control Design).\n\nBoth of these tuning methods assume a single-loop control configuration. Simulink Control Design software includes other tuning approaches that suit more complex configurations. For information about other ways to tune a PID Controller block, see Choose a Control Design Approach (Simulink Control Design).\n\n#### Dependencies\n\nTo enable this parameter, in the Main tab, set the controller-parameters Source to `internal`.\n\nZero-crossing detection can accurately locate signal discontinuities without resorting to excessively small time steps that can lead to lengthy simulation times. If you select Limit output or activate External reset in your PID Controller block, activating zero-crossing detection can reduce computation time in your simulation. Selecting this parameter activates zero-crossing detection:\n\n• At initial-state reset\n\n• When entering an upper or lower saturation state\n\n• When leaving an upper or lower saturation state\n\n#### Programmatic Use\n\n Block Parameter: `ZeroCross` Type: string, character vector Values: `\"on\"`, `\"off\"` Default: `\"on\"`\n\n### Initialization\n\nSimulink uses initial conditions to initialize the integrator and derivative-filter (or the unfiltered derivative) output at the start of a simulation or at a specified trigger event. (See the External reset parameter.) These initial conditions determine the initial block output. Use this parameter to select how to supply the initial condition values to the block.\n\n`internal`\n\nSpecify the initial conditions using the Integrator Initial condition and Filter Initial condition parameters. If Use filtered derivative is not selected, use the Differentiator parameter to specify the initial condition for the unfiltered differentiator instead of a filter initial condition.\n\n`external`\n\nSpecify the initial conditions externally using block inputs. Additional input ports Io and Do appear on the block. If Use filtered derivative is not selected, supply the initial condition for the unfiltered differentiator at Do instead of a filter initial condition.\n\n#### Programmatic Use\n\n Block Parameter: `InitialConditionSource` Type: string, character vector Values: `\"internal\"`, `\"external\"` Default: `\"internal\"`\n\nSimulink uses the integrator initial condition to initialize the integrator at the start of a simulation or at a specified trigger event (see External reset). The integrator initial condition and the filter initial condition determine the initial output of the PID controller block.\n\nThe integrator initial condition cannot be `NaN` or `Inf`.\n\n#### Dependencies\n\nTo use this parameter, in the Initialization tab, set Source to `internal`, and set Controller to a type that has integral action.\n\n#### Programmatic Use\n\n Block Parameter: `InitialConditionForIntegrator` Type: scalar, vector Default: 0\n\nSimulink uses the filter initial condition to initialize the derivative filter at the start of a simulation or at a specified trigger event (see External reset). The integrator initial condition and the filter initial condition determine the initial output of the PID controller block.\n\nThe filter initial condition cannot be `NaN` or `Inf`.\n\n#### Dependencies\n\nTo use this parameter, in the Initialization tab, set Source to `internal`, and use a controller that has a derivative filter.\n\n#### Programmatic Use\n\n Block Parameter: `InitialConditionForFilter` Type: scalar, vector Default: 0\n\nWhen you use an unfiltered derivative, Simulink uses this parameter to initialize the differentiator at the start of a simulation or at a specified trigger event (see External reset). The integrator initial condition and the derivative initial condition determine the initial output of the PID controller block.\n\nThe derivative initial condition cannot be `NaN` or `Inf`.\n\n#### Dependencies\n\nTo use this parameter, set Time domain to `Discrete-time`, clear the Use filtered derivative check box, and in the Initialization tab, set Source to `internal`.\n\n#### Programmatic Use\n\n Block Parameter: `DifferentiatorICPrevScaledInput` Type: scalar, vector Default: 0\n\nUse this parameter to specify whether to apply the Integrator Initial condition and Filter Initial condition parameter to the corresponding block state or output. You can change this parameter at the command line only, using `set_param` to set the `InitialConditionSetting` parameter of the block.\n\n`State (most efficient)`\n\nUse this option in all situations except when the block is in a triggered subsystem or a function-call subsystem and simplified initialization mode is enabled.\n\n`Output`\n\nUse this option when the block is in a triggered subsystem or a function-call subsystem and simplified initialization mode is enabled.\n\nThis parameter is only accessible through programmatic use.\n\n#### Programmatic Use\n\n Block Parameter: `InitialConditionSetting` Type: string, character vector Values: `\"state\"`, `\"output\"` Default: `\"state\"`\n\nSpecify the trigger condition that causes the block to reset the integrator and filter to initial conditions. (If Use filtered derivative is not selected, the trigger resets the integrator and differentiator to initial conditions.) Selecting any option other than `none` enables the Reset port on the block for the external reset signal.\n\n`none`\n\nThe integrator and filter (or differentiator) outputs are set to initial conditions at the beginning of simulation, and are not reset during simulation.\n\n`rising`\n\nReset the outputs when the reset signal has a rising edge.\n\n`falling`\n\nReset the outputs when the reset signal has a falling edge.\n\n`either`\n\nReset the outputs when the reset signal either rises or falls.\n\n`level`\n\nReset the outputs when the reset signal either:\n\n• Is nonzero at the current time step\n\n• Changes from nonzero at the previous time step to zero at the current time step\n\nThis option holds the outputs to the initial conditions while the reset signal is nonzero.\n\n#### Dependencies\n\nTo enable this parameter, set Controller to a type that has derivative or integral action.\n\n#### Programmatic Use\n\n Block Parameter: `ExternalReset` Type: string, character vector Values: `\"none\"`, `\"rising\"`, `\"falling\"`, `\"either\"`,`\"level\"` Default: `\"none\"`\n\nSelect to force Simulink and Simulink Control Design linearization commands to ignore any reset mechanism specified in the External reset parameter. Ignoring reset states allows you to linearize a model around an operating point even if that operating point causes the block to reset.\n\n#### Programmatic Use\n\n Block Parameter: `IgnoreLimit` Type: string, character vector Values: `\"off\"`, `\"on\"` Default: `\"off\"`\n\nSignal tracking lets the block output follow a tracking signal that you provide at the TR port. When signal tracking is active, the difference between the tracking signal and the block output is fed back to the integrator input with a gain `Kt`, specified by the Tracking gain (Kt) parameter. Signal tracking has several applications, including bumpless control transfer and avoiding windup in multiloop control structures.\n\n#### Bumpless control transfer\n\nUse signal tracking to achieve bumpless control transfer in systems that switch between two controllers. Suppose you want to transfer control between a PID controller and another controller. To do so, connecting the controller output to the TR input as shown in the following illustration.", null, "For more information, see Bumpless Control Transfer with a Two-Degree-of-Freedom PID Controller.\n\n#### Multiloop control\n\nUse signal tracking to prevent block windup in multiloop control approaches. For an example illustrating this approach with a 1DOF PID controller, see Prevent Block Windup in Multiloop Control.\n\n#### Dependencies\n\nTo enable this parameter, set Controller to a type that has integral action.\n\n#### Programmatic Use\n\n Block Parameter: `TrackingMode` Type: string, character vector Values: `\"off\"`, `\"on\"` Default: `\"off\"`\n\nWhen you select Enable tracking mode, the difference between the signal TR and the block output is fed back to the integrator input with a gain `Kt`. Use this parameter to specify the gain in that feedback loop.\n\n#### Dependencies\n\nTo enable this parameter, select Enable tracking mode.\n\n#### Programmatic Use\n\n Block Parameter: `Kt` Type: scalar Default: 1\n\n### Output saturation\n\nActivating this option limits the block output internally to the block, so that you do not need a separate Saturation block after the controller. It also allows you to activate the anti-windup mechanism built into the block (see the Anti-windup method parameter). Specify the saturation limits using the Lower saturation limit and Upper saturation limit parameters.\n\n#### Programmatic Use\n\n Block Parameter: `LimitOutput` Type: string, character vector Values: `\"off\"`, `\"on\"` Default: `\"off\"`\n\nSpecify the upper limit for the block output. The block output is held at the Upper saturation limit whenever the weighted sum of the proportional, integral, and derivative actions exceeds that value.\n\n#### Dependencies\n\nTo enable this parameter, select Limit output.\n\n#### Programmatic Use\n\n Block Parameter: `UpperSaturationLimit` Type: scalar Default: `Inf`\n\nSpecify the lower limit for the block output. The block output is held at the Lower saturation limit whenever the weighted sum of the proportional, integral, and derivative actions goes below that value.\n\n#### Dependencies\n\nTo enable this parameter, select Limit output.\n\n#### Programmatic Use\n\n Block Parameter: `LowerSaturationLimit` Type: scalar Default: `-Inf`\n\nForce Simulink and Simulink Control Design linearization commands to ignore block output limits specified in the Upper limit and Lower limit parameters. Ignoring output limits allows you to linearize a model around an operating point even if that operating point causes the block to exceed the output limits.\n\n#### Dependencies\n\nTo enable this parameter, select the Limit output parameter.\n\n#### Programmatic Use\n\n Block Parameter: `LinearizeAsGain` Type: string, character vector Values: `\"off\"`, `\"on\"` Default: `\"off\"`\n\nWhen you select Limit output and the weighted sum of the controller components exceeds the specified output limits, the block output holds at the specified limit. However, the integrator output can continue to grow (integrator windup), increasing the difference between the block output and the sum of the block components. In other words, the internal signals in the block can be unbounded even if the output appears bounded by saturation limits. Without a mechanism to prevent integrator windup, two results are possible:\n\n• If the sign of the signal entering the integrator never changes, the integrator continues to integrate until it overflows. The overflow value is the maximum or minimum value for the data type of the integrator output.\n\n• If the sign of the signal entering the integrator changes once the weighted sum has grown beyond the output limits, it can take a long time to unwind the integrator and return the weighted sum within the block saturation limit.\n\nIn either case, controller performance can suffer. To combat the effects of windup without an anti-windup mechanism, it may be necessary to detune the controller (for example, by reducing the controller gains), resulting in a sluggish controller. To avoid this problem, activate an anti-windup mechanism using this parameter.\n\n`none`\n\nDo not use an anti-windup mechanism.\n\n`back-calculation`\n\nUnwind the integrator when the block output saturates by feeding back to the integrator the difference between the saturated and unsaturated control signal. The following diagram represents the back-calculation feedback circuit for a continuous-time controller. To see the actual feedback circuit for your controller configuration, right-click on the block and select Mask > Look Under Mask.", null, "Use the Back-calculation coefficient (Kb) parameter to specify the gain of the anti-windup feedback circuit. It is usually satisfactory to set `Kb = I`, or for controllers with derivative action, `Kb = sqrt(I*D)`. Back-calculation can be effective for plants with relatively large dead time .\n\n`clamping`\n\nIntegration stops when the sum of the block components exceeds the output limits and the integrator output and block input have the same sign. Integration resumes when the sum of the block components exceeds the output limits and the integrator output and block input have opposite sign. Clamping is sometimes referred to as conditional integration.\n\nClamping can be useful for plants with relatively small dead times, but can yield a poor transient response for large dead times .\n\n#### Dependencies\n\nTo enable this parameter, select the Limit output parameter.\n\n#### Programmatic Use\n\n Block Parameter: `AntiWindupMode` Type: string, character vector Values: `\"none\"`, `\"back-calculation\"`,`\"clamping\"` Default: `\"none\"`\n\nThe `back-calculation` anti-windup method unwinds the integrator when the block output saturates. It does so by feeding back to the integrator the difference between the saturated and unsaturated control signal. Use the Back-calculation coefficient (Kb) parameter to specify the gain of the anti-windup feedback circuit. For more information, see the Anti-windup method parameter.\n\n#### Dependencies\n\nTo enable this parameter, select the Limit output parameter, and set the Anti-windup method parameter to `back-calculation`.\n\n#### Programmatic Use\n\n Block Parameter: `Kb` Type: scalar Default: 1\n\n### Data Types\n\nThe parameters in this tab are primarily of use in fixed-point code generation using Fixed-Point Designer™. They define how numeric quantities associated with the block are stored and processed when you generate code.\n\nIf you need to configure data types for fixed-point code generation, click Open Fixed-Point Tool and use that tool to configure the rest of the parameters in the tab. For information about using Fixed-Point Tool, see Autoscaling Data Objects Using the Fixed-Point Tool (Fixed-Point Designer).\n\nAfter you use Fixed-Point Tool, you can use the parameters in this tab to make adjustments to fixed-point data-type settings if necessary. For each quantity associated with the block, you can specify:\n\n• Floating-point or fixed-point data type, including whether the data type is inherited from upstream values in the block.\n\n• The minimum and maximum values for the quantity, which determine how the quantity is scaled for fixed-point representation.\n\nFor assistance in selecting appropriate values, click", null, "to open the Data Type Assistant for the corresponding quantity. For more information, see Specify Data Types Using Data Type Assistant.", null, "The specific quantities listed in the Data Types tab vary depending on how you configure the PID controller block. In general, you can configure data types for the following types of quantities:\n\n• Product output — Stores the result of a multiplication carried out under the block mask. For example, P product output stores the output of the gain block that multiplies the block input with the proportional gain P.\n\n• Parameter — Stores the value of a numeric block parameter, such as P, I, or D.\n\n• Block output — Stores the output of a block that resides under the PID controller block mask. For example, use Integrator output to specify the data type of the output of the block called Integrator. This block resides under the mask in the Integrator subsystem, and computes integrator term of the controller action.\n\n• Accumulator — Stores values associated with a sum block. For example, SumI2 Accumulator sets the data type of the accumulator associated with the sum block SumI2. This block resides under the mask in the Back Calculation subsystem of the Anti-Windup subsystem.\n\nIn general, you can find the block associated with any listed parameter by looking under the PID Controller block mask and examining its subsystems. You can also use the Model Explorer to search under the mask for the listed parameter name, such as `SumI2`. (See Model Explorer.)\n\nMatching Input and Internal Data Types\n\nBy default, all data types in the block are set to `Inherit: Inherit via internal rule`. With this setting, Simulink chooses data types to balance numerical accuracy, performance, and generated code size, while accounting for the properties of the embedded target hardware.\n\nUnder some conditions, incompatibility can occur between data types within the block. For instance, in continuous time, the Integrator block under the mask can accept only signals of type `double`. If the block input signal is a type that cannot be converted to `double`, such as `uint16`, the internal rules for type inheritance generate an error when you generate code.\n\nTo avoid such errors, you can use the Data Types settings to force a data type conversion. For instance, you can explicitly set P product output, I product output, and D product output to `double`, ensuring that the signals reaching the continuous-time integrators are of type `double`.\n\nIn general, it is not recommended to use the block in continuous time for code generation applications. However, similar data type errors can occur in discrete time, if you explicitly set some values to data types that are incompatible with downstream signal constraints within the block. In such cases, use the Data Types settings to ensure that all data types are internally compatible.\n\n#### Fixed-Point Operational Parameters\n\nSpecify the rounding mode for fixed-point operations. For more information, see Rounding (Fixed-Point Designer).\n\nBlock parameters always round to the nearest representable value. To control the rounding of a block parameter, enter an expression using a MATLAB® rounding function into the mask field.\n\n#### Programmatic Use\n\n Block Parameter: `RndMeth` Type: character vector Values: ```'Ceiling' | 'Convergent' | 'Floor' | 'Nearest' | 'Round' | 'Simplest' | 'Zero'``` Default: `'Floor'`\n\nSpecify whether overflows saturate or wrap.\n\n• `off` — Overflows wrap to the appropriate value that the data type can represent.\n\nFor example, the number 130 does not fit in a signed 8-bit integer and wraps to -126.\n\n• `on` — Overflows saturate to either the minimum or maximum value that the data type can represent.\n\nFor example, an overflow associated with a signed 8-bit integer can saturate to -128 or 127.\n\n### Tip\n\n• Consider selecting this check box when your model has a possible overflow and you want explicit saturation protection in the generated code.\n\n• Consider clearing this check box when you want to optimize efficiency of your generated code.\n\nClearing this check box also helps you to avoid overspecifying how a block handles out-of-range signals. For more information, see Troubleshoot Signal Range Errors.\n\n• When you select this check box, saturation applies to every internal operation on the block, not just the output or result.\n\n• In general, the code generation process can detect when overflow is not possible. In this case, the code generator does not produce saturation code.\n\n#### Programmatic Use\n\n Block Parameter: `SaturateOnIntegerOverflow` Type: character vector Values: `'off' | 'on'` Default: `'off'`\n\nSelect this parameter to prevent the fixed-point tools from overriding the data types you specify on this block. For more information, see Lock the Output Data Type Setting (Fixed-Point Designer).\n\n#### Programmatic Use\n\n Block Parameter: `LockScale` Type: character vector Values: `'off' | 'on'` Default: `'off'`\n\n### State Attributes\n\nThe parameters in this tab are primarily of use in code generation.\n\nAssign a unique name to the state associated with the integrator or the filter, for continuous-time PID controllers. (For information about state names in a discrete-time PID controller, see the State name parameter.) The state name is used, for example:\n\n• For the corresponding variable in generated code\n\n• As part of the storage name when logging states during simulation\n\n• For the corresponding state in a linear model obtain by linearizing the block\n\nA valid state name begins with an alphabetic or underscore character, followed by alphanumeric or underscore characters.\n\n#### Dependencies\n\nTo enable this parameter, set Time domain to `Continuous-time`.\n\n#### Programmatic Use\n\n Parameter: `IntegratorContinuousStateAttributes`, `FilterContinuousStateAttributes` Type: character vector Default: `''`\n\nAssign a unique name to the state associated with the integrator or the filter, for discrete-time PID controllers. (For information about state names in a continuous-time PID controller, see the State name (e.g., 'position') parameter.)\n\nA valid state name begins with an alphabetic or underscore character, followed by alphanumeric or underscore characters. The state name is used, for example:\n\n• For the corresponding variable in generated code\n\n• As part of the storage name when logging states during simulation\n\n• For the corresponding state in a linear model obtain by linearizing the block\n\nFor more information about the use of state names in code generation, see Apply Storage Classes to Individual Signal, State, and Parameter Data Elements (Simulink Coder).\n\n#### Dependencies\n\nTo enable this parameter, set Time domain to `Discrete-time`.\n\n#### Programmatic Use\n\n Parameter: `IntegratorStateIdentifier`, `FilterStateIdentifier` Type: string, character vector Default: `\"\"`\n\nSelect this parameter to require that the discrete-time integrator or filter state name resolves to a Simulink signal object.\n\n#### Dependencies\n\nTo enable this parameter for the discrete-time integrator or filter state:\n\n1. Set Time domain to `Discrete-time`.\n\n2. Specify a value for the integrator or filter State name.\n\n3. Set the model configuration parameter Signal resolution to a value other than `None`.\n\nSelecting this check box disables Code generation storage class for the corresponding integrator or filter state.\n\n#### Programmatic Use\n\n Block Parameter: `IntegratorStateMustResolveToSignalObject`, `FilterStateMustResolveToSignalObject` Type: string, character vector Values: `\"off\"`, `\"on\"` Default: `\"off\"`\n\nSelect state storage class for code generation. If you do not need to interface to external code, select `Auto`.\n\nFor more information, see Apply Storage Classes to Individual Signal, State, and Parameter Data Elements (Simulink Coder) and Apply Built-In and Customized Storage Classes to Data Elements (Embedded Coder).\n\n#### Dependencies\n\nTo enable this parameter for the discrete-time integrator or filter state:\n\n1. Set Time domain to `Discrete-time`.\n\n2. Specify a value for the integrator or filter State name.\n\n3. Set the model configuration parameter Signal resolution to a value other than `None`.\n\n#### Programmatic Use\n\n Block Parameter: `IntegratorRTWStateStorageClass`, `FilterRTWStateStorageClass` Type: string, character vector Values: `\"Auto\"`, `\"ExportedGlobal\"`, `\"ImportedExtern\"` | `\"ImportedExternPointer\"` Default: `\"Auto\"`\n\nSpecify a storage type qualifier such as `const` or `volatile`.\n\n### Note\n\nThis parameter will be removed in a future release. To apply storage type qualifiers to data, use custom storage classes and memory sections. Unless you use an ERT-based code generation target with Embedded Coder®, custom storage classes and memory sections do not affect the generated code.\n\n#### Dependencies\n\nTo enable this parameter, set Code generation storage class to any value other than `Auto`.\n\n#### Programmatic Use\n\n Block Parameter: `IntegratorRTWStateStorageTypeQualifier`,`FilterRTWStateStorageTypeQualifier` Type: string, character vector Values:`\"\"`,`\"const\"`,`\"volatile\"` Default: `\"\"`\n\n## Block Characteristics\n\n Data Types `double` | `fixed point` | `integer` | `single` Direct Feedthrough `no` Multidimensional Signals `no` Variable-Size Signals `no` Zero-Crossing Detection `no`" ]
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http://www.scitlion.com/index.php?m=content&c=index&a=show&catid=151&id=230
[ "## Contact Us\n\n+86-10-81508344(Sales Department)\nNo.5, Xingguang 4 Ave ,Opto-Mechatronics Industrial Park ,TongZhou district,Beijing 101111,China\n\n## Nonlinear Crystals", null, "• Product Name: LiNbO3\n• Dimension Tolerance:W x H (+/- 0.1mm) x L (+0.2/-0.1mm)\n• Angle Tolerance:Δθ ≤ 0.25°, Δφ ≤ 0.25°\n• Parallelism:< 20 arc seconds\n• Perpendicularity:< 5 arc minutes\n• Chamfer:0.1mm at 45°\n• Clear Aperture:90%\n• Scratch and Dig:20/10 or better\n• Surface Flatness:λ/8 at 633nm\n• Wavefront Distortion:λ/8 at 633nm\n• Damage Threshold :100 MW/cm2 at 1064nm, 10ns, 10Hz\n\n## Product overview\n\nLithium Niobate (LiNbO3) or LN Crystal is widely used as frequency doublers for wavelength over 1μm and optical parametric oscillators (OPOs) pumped at 1064 nm and quasi-phase-matched (QPM) devices. Additionally due to its large Electro-Optic (EO) and Acousto-Optic (AO) coefficients, LiNbO3 is the most commonly used material for Pockels Cell, Q-switches and phase nodulators, waveguide substrate, and surface acoustic wave (SAW) wafers, etc.\n\nScitlion is able to offer\nSize upto 10 x 10 x 20mm\nLength from 0.02mm to 20mm\nCoating upon your request", null, "", null, "Density 4.64 g/cm3 Melting Point 1253°C Curie Temperature 1140°C Mohs Hardness 5 Elastic Stiffness Coefficients CE11= 2.33(×1011N/m2  ) CE33= 2.77(×1011N/m2  ) Crystal Structure Trigonal Lattices Constants a=5.148A , c=13.863A Thermal and Electrical Properties Thermal Expansion Coefficients 2.0x10-6/K  (// a-axis  at 25°C) 2.2x10-6/K  (// c-axis  at 25°C) Thermal Conductivity 38W/m/K  at 25°C Resistivity 2x10-6Ω·cm@200°C Piezoelectric Strain Constant D22=2.04(x10-11C/N) D33=19.22(x10-11C/N) Optical and Electro-optical Properties Transparency Range 420nm to 5200nm Optical Homogeneity ~ 5 x 10-5 /cm Refractive indices ne = 2.146, no = 2.220  at 1300 nm ne = 2.156, no = 2.232  at 1064 nm ne = 2.203, no = 2.286  at 632.8 nm NLO Coefficients: d33 = 86 x d36 (KDP) d31 =11.6 x d36 (KDP) d22 = 5.6 x d36 (KDP) Effective NLO Coefficients: deff(I)=d31sinθ-d22cossin3Φ deff(II)=d22cos2θcos3Φ Electro-Optic Coefficients r T33 = 32 pm/V,   rS33 = 31 pm/V, rT31 =10 pm/V,     rS31=8.6 pm/V, rT22 = 6.8 pm/V,   rS22= 3.4 pm/V, Half-Wave Voltage 3.03 KV  (Electrical field // z, light ┴ z) 4.02 KV  (Electrical field // x or y, light // z) Sellmeier Equations (λ in μm) no2=4.9048+0.11768/(λ2-0.04750)-0.027169λ2 ne2=4.5820+0.099169/(λ2-0.04443)-0.02195λ2" ]
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https://asp-eurasipjournals.springeropen.com/articles/10.1186/s13634-021-00730-w
[ "# Multi-source and multi-fault condition monitoring based on parallel factor analysis and sequential probability ratio test\n\n## Abstract\n\nThe monitoring of mechanical equipment systems contains an increasing number of complex content, expanding from traditional time, and frequency information to three-dimensional data of the time, space, and frequency information, and even higher-dimensional data containing subjects, experimental conditions. For high-dimensional data analysis, traditional decomposition methods such as Hilbert transform, fast Fourier transformation, and Gabor transformation not only lose the integrity of the data, but also increase the amount of calculation and introduce a lot of redundant information. The phenomenon of feature coupling, aliasing, and redundancy between the mechanical multi-source data signals will cause the inaccuracy of the evaluation, diagnosis, and prediction of industrial production operation status. The analysis of the three-way tensor composed of channel, frequency, and time is called parallel factor analysis (PARAFAC). The properties between the parallel factor analysis results and the input signals are studied through simulation experiments. Parallel factor analysis is used to decompose the third-order tensor composed of channel-time-frequency after continuous wavelet transformation of vibration signal into channel, time, and frequency characteristics. Multi-scale parallel factor analysis successfully extracted non-linear multi-dimensional dynamic fault characteristics by generating the spatial, spectral, time-domain signal loading value and three-dimensional fault characteristic expression. In order to verify the effectiveness of the space, frequency, and time domain signal loading values of the fault characteristic factors generated by the centrifugal pump system after parallel factor analysis, the characteristic factors obtained after parallel factor analysis are used as the SPRT test sequence for identification and verification. The results indicate that the method proposed in this article improves the measurement accuracy and intelligence of mechanical fault detection.\n\n## 1 Introduction\n\nIn recent years, equipment fault diagnosis, as a new technology that crosses various disciplines, has been developed rapidly and has produced huge economic benefits [1,2,3,4,5,6]. The centrifugal pump is an important energy conversion and liquid transmission device in the process industry and its working state directly affects the production of the entire operating equipment. A slight damage to the impeller will shorten the running time of the centrifugal pump and disturb the operation of the equipment. When the impeller fails, it will cause damage to the centrifugal pump components or personal injury accidents, which will cause significant economic losses [7, 8]. The normal operation and failure of the centrifugal pump will cause the equipment to vibrate. The vibration signal contains rich information of the pump body running state and is easy to be collected, which can be used to monitor and diagnose the running state of the centrifugal pump. Fault diagnosis is generally divided into three steps: first, we collect the relevant vibration signal of the diagnostic object; then, the signal is analyzed and processed to acquire the characteristics of the vibration signal; finally, pattern recognition and fault diagnosis are performed through the corresponding extracted special diagnosis [9, 10]. The core content is to obtain the effective characteristics of the vibration signal. Due to the complex structure of the centrifugal pump, many excitation sources and mutual interference, the vibration signal of the centrifugal pump is a non-linear and non-stationary signal. Researchers have proposed various effective diagnostic methods to process the collected raw vibration signals of the centrifugal pump, extract effective information, and improve the accuracy of diagnosis. Wenjian Huang et al. extracted the characteristic parameters of the vibration signal through time-domain signal analysis, then the PCA was used to reduce the amount of data, and finally the main component with the largest contribution rate was used as the input signal of SPRT to evaluate the proposed algorithm. Literature proposed an improved deep convolutional neural network (CNN) to identify defects in centrifugal pumps by using sound and image recognition. A feature extraction method based on empirical mode decomposition (EMD) was developed to detect the gravity of cavitation in the centrifugal pump by Azizi et al. . Liu Yang et al. proposed the new method for analysis of big data based on particle swarm optimization wavelet neural network for diagnosis in the gearbox. Literature applies variational mode decomposition (VMD) with different input parameters to fault diagnosis of multi-stage pumps. Signal processing combined with empirical mode decomposition (EMD) and fuzzy c-means clustering is used for monitoring piston pump defects in literature . The traditional decomposition method of processing high-dimensional data will not only lose the integrity of the data but also increase the amount of calculation and introduce redundancy [15,16,17,18]. These methods of extracting time-frequency characteristics from single-channel signals, such as Fourier transform, cannot reflect the internal relationship of non-linear changes between multi-source channel characteristic signals, nor can they eliminate information interference.\n\nMechanical non-linear multi-fault mode multi-source dynamic feature identification is a technical bottleneck and difficult problem encountered in the application of fault diagnosis in process industry production lines. It not only needs to extract the time-frequency characteristics of multi-source fault signals, but also to ensure the correspondence between non-linear variables and multi-fault modes and multi-source fault features in time, frequency, and space after feature extraction. Parallel factor analysis proposed by Carroll, Chang , and Harshman in 1970 is a three-dimensional or multi-dimensional signal processing algorithm that uses iterative least squares to resolve the decomposition and identification of multi-dimensional matrices. The general time-frequency decomposition ignores the spatial information of the vibration signal and cannot handle multi-channel data [21,22,23]. Data framing in the form of a three-way array indexed by channel, frequency and time allows the application of a unique decomposition called Parallel Factor Analysis (PARAFAC). The decomposition uniqueness of PARAFAC model can obtain its model parameters without ambiguity so that the PARAFAC model has important application value. As a data processing algorithm, PARAFAC model has been successfully applied in fluorescence spectroscopy, psychology, signal processing, food science, and other fields. Multi-channel electroencephalogram EEG data can usually be expressed as an M×N×P three-way data set and the components of the three-way data array correspond to the channel (electrodes at different positions), time (data samples) and frequency components. Schmitz, S. applied PARAFAC to analyze the temporal and spatial patterns of functional connections between neurons, which were revealed in the sequence of peaks recorded in the cat’s main visual cortex (area 18) . This parallel factor analysis was applied for decomposing EEG data into space-time-frequency components during the resting state and mental arithmetic by Miwakeichi, Fumikazu et al. . Rost'akova compares non-negative Tucker decomposition with parallel factor analysis to identify and measure human brain electrical rhythms . In the literature , Choi, Ji Yeh proposed a new extension function PARAFAC for processing response to three-dimensional data arranged along a two-dimensional domain and one-dimensional parameters. Technically, this method combines PARAFAC with basis function expansion approximation and is applied to EEG data to prove its empirical usefulness. A parallel factor analysis study showed that the frontal lobe area with higher frequency response is the main feature of laser evoked potential in rats with chronic inflammatory pain .\n\nParallel decomposition has attracted great attention, because parallel factorization can process the constructed high-dimensional data as a whole, which not only retains the overall structure information of the data, but also makes the structure more compact and easy to understand. In the literature , parallel factor analysis was used as the diagnostic tool through decomposing centrifugal pump diagnostic signal into time-frequency-space modes. Considering the difficulty of extracting fault features from rolling bearings under strong background noise, Yang Cheng proposed a new method based on variable mode decomposition (VMD) and phase space parallel factor analysis to detect weak fault signals of rolling bearings. In order to overcome the inability to extract sparse and interpretable latent variables from batch data, literature proposed a batch three-way data array sparse model based on sparse parallel factor (SPARAFAC) decomposition. Sparse factor matrices have the potential advantage of being easy to interpret because they eliminate redundant data information and show significant variable correlation. In chemistry, medicine, and food science commonly used fluorescence excitation and emission data typically contain several chemical components at different concentrations. Fluorescence spectroscopy can generate a three-way data set with the mode “sample × excitation × emission.” The main purpose of the analysis of this data type is to determine which chemicals are present in each sample and their relative concentrations. Reference conducted a comparison between parallel factor analysis (PARAFAC) and support vector machine (SVM) to identify and distinguish the fluorescence spectrum of coconut water brands. The above results indicate that fluorescence spectroscopy combined with PARAFAC and SVM method has been proved to be a simple and rapid detection method for coconut water and other beverages. This study aims to determine whether the composition or distribution of humus in lake sediments can be characterized by chemometric spectral data. This method determines the three-dimensional excitation emission matrix in the extracted humus and performs spectral analysis of the data by using parallel factor analysis (PARAFAC) with classification tree and regression tree (CART).\n\nThe theory of sequential probability ratio test (SPRT) that is a branch of mathematical statistics was proposed by Abraham Wald in 1947 in order to solve the problem of sampling and acceptance of valuable military products. This method provides an approximate formula for the critical value of accepting the null hypothesis H0 or accepting the alternative hypothesis H1 based on the sample values obtained from each observation, and also provides the average sampling times and power function of this test method. In 1948, Abraham Wald and American statistician Wolfowitz proved that the above-mentioned sequential probability is the smallest number of sampling times required for the test in all the two types of tests whose error probability does not exceed α and β, respectively. The sequential probability ratio test is the most fundamental sequential test in sequential analysis proposed by Abraham Wald, and it has subsequently been widely developed in various fields. Almost all the hypothesis testing problems of SPRT in mechanical fault diagnosis, such as signal detection, model variable point detection, life data analysis of centrifugal pump, and crack detection of gearbox, can be well applied. This research were performed on actual faults in a laboratory-scale distillation plan based on artificial neural network-multi-layer perceptron (ANN-MLP) and the Wald sequential probability ratio test (SPRT). In the literature , Guo Peng proposed Gaussian process and SPRT wind turbine power curve modeling and monitoring. The modeling and monitoring method proposed in this paper successfully identified two wind anemometer failures and pitch system failures. Literature proposed a fault detection algorithm based on the sequential probability ratio test (SPRT) and chi-square test for redundant multi-sensor navigation systems for supersonic cruise ships (HCV).\n\nThe rest of this paper is organized as follows. In Section 2, the parallel factorization model and simulation is described. The multi-scale parallel factorization optimization algorithm for non-linear multi-source fault characteristic signal extraction is established. The characteristic factor signal is successfully obtained from the matrix factor, and the “loading” factor and “component” factor are defined. In Section 3, we studied the multi-channel data decision theory based on SPRT, and established an adaptive optimization diagnosis method for tracking and identifying the non-linear multi-dimensional dynamic optimal characteristic signal. To research the validity of the multi-scale parallel factor analysis and SPRT for multi-channel signal in actual complex industrial production, the centrifugal pump fault diagnosis experimental system was designed and implemented in Section 4. Following that, multi-source dynamic feature extraction based on parallel factorization and SPRT for the multi-source condition monitoring of centrifugal pump are presented in Section 5. Finally, the conclusions are drawn in Section 6.\n\n## 2 The model and simulation\n\n### 2.1 Parallel factorization model\n\nIn a two-dimensional matrix, xi, j is generally used to represent the elements in the two-dimensional matrix (i represents any row, j represents any column). Similarly, in the three-dimensional matrix, we use xi, j, kto represent any element in the three-dimensional matrix. At present, there is no definite name naming subscript k, let us call it “page” . The subscript of the three-dimensional matrix consists of three index value row, column, and page composition. The left picture in Fig. 1 shows the three-dimensional matrix and the right picture shows its sub-matrix. When a certain dimension in a three-dimensional matrix is fixed, it constitutes a sub-matrix of the three-dimensional matrix that is called a slice of the three-dimensional matrix along a certain dimension.\n\nThe expansion of the three-dimensional matrix is actually to rearrange the slices of the three-dimensional matrix to constitute the new two-dimensional matrix. For example, we fixed the rows and columns of a three-dimensional matrix and rearranged its pages to formulate the new two-dimensional matrix. At this time, the number of rows is equivalent to the number of rows I of the original matrix and the number of columns changes from the original J to J × K, denoted asXI × JK. It is expressed as shown in Eq. (1).\n\n$${X}^{I\\times JK}=\\left[{X}_{k=1},{X}_{k=2}\\cdots {X}_{k=K}\\right]$$\n(1)\n\nOf course, it can also be expanded by column, such as XIK × J, which is defined as formula (2).\n\n$${X}^{IK\\times J}=\\left[\\begin{array}{c}{X}_{k=1}\\\\ {}{X}_{k=2}\\\\ {}\\vdots \\\\ {}{X}_{k=K}\\end{array}\\right]$$\n(2)\n\nAfter expanding by columns, we acknowledge that I × K is displayed as the number of rows of the new matrix and parameter “J” is the number of columns.\n\nThe symbol $${x}_{i,j,k}=\\sum \\limits_{f=1}^F{a}_{i,f}{b}_{j,f}{c}_{k,f}$$ can be used to express any element in a three-dimensional or larger than three-dimensional matrix, the variables i, j, and k in the formula can be any natural numbers. The elements in the i-th row, j-th column, and k-th page of the matrix X can be represented by xi, j, k. According to the definition, the low-rank decomposition of a two-dimensional matrix can be popularized to the low-rank decomposition of a three-dimensional matrix. Let the element XCI × J × Kof the three-dimensional matrix be defined as xijk, the variables i, j, and k in the formula can be any natural numbers. Similarly, it can be seen that the three-dimensional matrix can be indicated as the modality of the vector outer product of the following formula (3).\n\n$$X={a}_1\\circ {b}_1\\circ {c}_1+,\\dots, +{a}_R\\circ {b}_R\\circ {c}_R=\\sum \\limits_{r=1}^R{a}_r\\circ {b}_r\\circ {c}_r$$\n(3)\n\nIt gives the process of a three-dimensional matrix low rank decomposition in formula (3) and the symbol R of the formula is indicated as the rank of the three-dimensional matrix X, where crCK, brCJ,arCI, r = 1,...,R. Harshman names the low-rank decomposition model of three-dimensional matrix given by formula (3) as the general model of parallel factor. The general forms of the parallel factorization model are interpreted in Fig. 2, in which X can be displayed by the popular geometric cube.\n\nTheseC = [c1, …, cR];B = [b1, …, bR];A = [a1, …, aR] are any given three two-dimensional matrix definitions. We define A, B, and C as the three loading matrices of the parallel factorization model of the general form. Equation (4) is shown as the scalar form of the general parallel factorization model. They are labeled as ckr = [C]k, r, bjr = [B]j, r,air = [A]i, r, and $${\\underline{x}}_{ijk}={\\left[\\underline{X}\\right]}_{i,j,k}$$.\n\n$${\\underline{x}}_{ijk}=\\sum \\limits_{r=1}^R{a}_{ir}{b}_{jr}{c}_{kr}$$\n(4)\n\nThe general form of the parallel factorization model can be viewed as the low-rank decomposition of a two-dimensional matrix extending to a three-dimensional matrix. The formula (4) indicates that these subitems $${\\underline{x}}_{ijk}$$of the three-dimensional matrix X can also be denoted as the sum of the products of R elements a, b, and c. Compared with the matrix elements xij in PCA, $${\\underline{x}}_{ijk}$$ has three independently changing dimensions called “mode A,” “mode B,” and “mode C.”\n\n### 2.2 Matrix essential equalization\n\nThere is a matrix ACI × J. If the matrix satisfies $$\\overline{A}=A\\Pi \\Delta$$, it is said that A and $$\\overline{A}$$ are matrix essential equalization, denoted as $$A\\cong \\overline{A}$$. Among them, ∏ is the column exchange matrix and is the diagonal scale matrix. There is one and only one non-zero element “1” in each row and each column of the column exchange matrix ∏. The function of the column exchange matrix is to rearrange the column vector of A in the order of ∏ without changing the value of the elements in the vector. The diagonal scale matrix is a J×J diagonal matrix with non-zero diagonal elements. The function of is to multiply each column of matrix A by a non-zero amplitude.\n\nAccording to the concept of matrix essential equalization, we take a two-dimensional matrixX = ABTas an example, where ACM × F,BCN × F. For any matrix $$\\overline{A}\\in {C}^{M\\times F},\\overline{B}\\in {C}^{N\\times F}$$, if $$X={AB}^T=\\overline{A}{\\overline{B}}^T$$is satisfied, then we can get to formula (5).\n\n$$\\overline{A}=A{\\prod}_1{\\Delta}_1,\\overline{B}=B{\\prod}_2{\\Delta}_2$$\n(5)\n\n∏2 and ∏2 in the formula are column exchange matrices, which means to rearrange the columns of the A and B matrices. 1 and 2 are diagonal scale matrices, which means that each column of matrix A and B is multiplied by the non-zero coefficient. At this time, the matrix decomposition is said to be unique.\n\nWhen the matrix factorization is unique, the matrix $$\\overline{A},\\overline{B}$$obtained by the matrix factorization is not completely equal to the original matrices A and B, they are only the essential equality relationship of the matrix. The essential equal relationship of matrix A and B is shown in the following formula (6) (7).\n\n$$\\overline{A}=A{\\prod}_A{\\Delta}_A\\cong A$$\n(6)\n$$\\overline{B}=B{\\prod}_B{\\Delta}_B\\cong B$$\n(7)\n\nDue to the existence of column exchange matrices ∏A, ∏B and diagonal scale matrices A, B, the order and magnitude of the column vectors in matrix $$\\overline{A},\\overline{B}$$ can be different from those of A and B. In matrix theory, they are used to be called column blur and scale blur, which are represented by column exchange matrix ∏ and diagonal scale matrix , respectively. In the process of matrix decomposition, if there is no structural constraint on the matrices A and B, column blur and scale blur are unavoidable. The above problem can be explained in the vector form of matrix decomposition, which is denotedA = [a1, , aF], B = [b1, , bF], whereafCI × 1bfCJ × 1 (f=1,...,F) is the column vector of A and B, respectively. The above formula can be expressed as the following vector form.\n\n$$X={AB}^T={\\mathrm{a}}_1{\\mathrm{b}}_1^T+{\\mathrm{a}}_2{\\mathrm{b}}_2^T+\\cdots +{\\mathrm{a}}_F{\\mathrm{b}}_F^T$$\n(8)\n\nIn formula (8), $${\\mathrm{a}}_1{\\mathrm{b}}_1^T,\\cdots, {\\mathrm{a}}_F{\\mathrm{b}}_F^T$$ are F matrices with rank 1. At this time, if the order of $${\\mathrm{a}}_1{\\mathrm{b}}_1^T$$,...,$${\\mathrm{a}}_F{\\mathrm{b}}_F^T$$ is changed arbitrarily, the value of matrix X is unchanged. Similarly, if the vector af is multiplied by the non-zero coefficient λf and the corresponding bf is multiplied by a non-zero coefficient 1/λf, the value of X will not change either. Assuming that the order of $${\\mathrm{a}}_1{\\mathrm{b}}_1^T$$and $${\\mathrm{a}}_2{\\mathrm{b}}_2^T$$ in formula (8) is exchanged, it can be rewritten into the following form:\n\n$$\\begin{array}{c}X={AB}^T={\\mathrm{a}}_1{\\mathrm{b}}_1^T+{\\mathrm{a}}_2{\\mathrm{b}}_2^T+\\cdots +{\\mathrm{a}}_F{\\mathrm{b}}_F^T\\\\ {}={\\lambda}_2{a}_2\\frac{1}{\\lambda_2}{b}_2^T+{\\lambda}_1{a}_1\\frac{1}{\\lambda_1}{b}_1^T+\\cdots +{\\lambda}_F{a}_F\\frac{1}{\\lambda_F}{b}_F^T\\\\ {}=A{\\prod}_A{\\Delta}_A{\\left(B{\\prod}_B{\\Delta}_B\\right)}^T\\\\ {}=\\overline{A}{\\overline{B}}^T\\end{array}}$$\n(9)\n\nThe result $$X={AB}^T=\\overline{A}{\\overline{B}}^T$$ can be obtained, the code in the equation is specifically expressed as follows $$\\overline{A}=A{\\prod}_A{\\Delta}_A$$,$$\\overline{B}=B{\\prod}_B{\\Delta}_B$$.\n\nIt can be seen that there are always $$\\overline{A}$$ and $$\\overline{B}$$ to achieve matrix decomposition, but they are only essentially equal to A and B. Therefore, column ambiguity and scale ambiguity are inherent ambiguities in the matrix decomposition process. Without additional constraints, the order and magnitude of loading matrix columns cannot be determined through matrix decomposition. Therefore, the uniqueness of matrix factorization given by the definition can also be called “essential uniqueness.” In the actual application process, some methods can be used to eliminate column blur and proportion blur caused by matrix decomposition.\n\n### 2.3 Recognizability and uniqueness of parallel factorization\n\nThe essential feature of the parallel factorization model is the uniqueness of the model. When there is no array blur, the matrices A, B, and C can be identified. The following conclusions can be obtained. When Xi = BDi(A)CT,i=1,2,...,I, ACI × F,BCJ × F,CCK × F is given, if kA + kB + kC ≥ 2F + 2, then these matrices A, B, and C are uniqueness for column exchange and plurality transformation or scale transformation.\n\nThe matrix composed of relatively independent columns taken from the absolute continuous distribution has full k-rank. If all three matrices meet this condition, the sufficient condition for recognizability is shown in formula (10).\n\n$$\\min \\left(I,F\\right)+\\min \\left(J,F\\right)+\\min \\left(K,F\\right)\\ge 2F+2$$\n(10)\n\nIf the matrices A, B, and C have other structural constraints, better identifiable results may be obtained. PARAFAC uniqueness theorem can be used to obtain the ith submatrix of the X-axis of PARAFAC model:\n\n$${X}_i^{J\\times K}={BD}_i(A){C}^T,i=1,\\dots, I$$\n(11)\n\nThe matrix in formula (11) satisfies the following ARI × R, BRJ × R, CRK × R. If the following conditions are met in Eq. (12).\n\n$${k}_A+{k}_B+{k}_C\\ge 2\\left(R+1\\right)$$\n(12)\n\nEven if there is column blur and scale blur, the matrix A, B, and C are unique. In mathematical language, when formula (13) is satisfied\n\n$${X}_i^{J\\times K}=\\hat{B}{D}_i\\left(\\hat{A}\\right){\\hat{C}}^T,i=1,\\dots, I$$\n(13)\n\nThe relation shown in formula (14) can be obtained.\n\n$$\\hat{A}=A\\prod {\\Delta}_1,\\hat{B}=B\\prod {\\Delta}_2,\\hat{C}=C\\prod {\\Delta}_3$$\n(14)\n\nEquation (14) shows that ∏ is a column fuzzy matrix, Δ1Δ2 and Δ3 are the scale fuzzy matrix and the equation of Δ1Δ2Δ3 = I needs to be satisfied.\n\n### 2.4 Trilinear alternating least square for parallel factor analysis\n\nThere are many methods to achieve the decomposition of PARAFAC, and the trilinear alternate least squares (TALS) algorithm is the most widely adopted methodology for data detection of parallel factor trilinear models. The fundamental principle of the TALS is to update the matrix in each step. First of all, TALS is updated by employing least squares (LS) to renovate the residual matrix based on the results of the previous estimate; then, it continues to update other matrices; finally, stop running until the result converges or reaches the set number of iterations after repeating the above steps. The trilinear model of three-dimensional data set X has the configuration shown in formula (15) below.\n\n$${x}_{i,j,k}=\\sum \\limits_{f=1}^F{a}_{i,f}{b}_{j,f}{c}_{k,f}+{e}_{ijk},i=1\\dots I;j=1\\dots J;k=1\\dots K$$\n(15)\n\nWhere F is the number of factors, ai,f is the i-th element in vector af, bj,f is the j-th element in vector bf, and ck,f is the k-th element in vector cf. The data set X of third-order tensor I × J × K is indicated as “xi,j,k.” The “eijk” represents the error set E of the third-order tensor I × J × K. Equation A = [a1, a2, , aI] is defined as the I × F matrix; the B = [b1, b2, , bJ], and C = [c1, c2, , cK] are defined as the J × F matrix and the K × F matrix.\n\n1. (1)\n\nFirst, matrix A is calculated by formula (16).\n\n$$\\left[\\begin{array}{l}{X}_{\\mathrm{K}1}\\\\ {}{X}_{\\mathrm{K}2}\\\\ {}\\mathrm{M}\\\\ {}{X}_{\\mathrm{K}\\ K}\\end{array}\\right]=\\left[\\begin{array}{l} BdiagC\\left(1,:\\right)\\\\ {} BdiagC\\left(2,:\\right)\\\\ {}\\mathrm{M}\\\\ {} BdiagC\\left(K,:\\right)\\end{array}\\right]{A}^T+{E}_K$$\n(16)\n\nFormula (16) satisfies XK k = Bdiag(C(k,:))AT +EK k. The error is expressed in terms of EK. The least squares (LS) estimate of AT is calculated by Eq. (17).\n\n$${A}^T={\\left[\\begin{array}{c} BdiagC\\left(1,:\\right)\\\\ {} BdiagC\\left(2,:\\right)\\\\ {}M\\\\ {} BdiagC\\left(2,:\\right)\\end{array}\\right]}^{+}\\left[\\begin{array}{c}{X}_{\\mathrm{K}1}\\\\ {}{X}_{\\mathrm{K}2}\\\\ {}\\mathrm{M}\\\\ {}{X}_{\\mathrm{K}K}\\end{array}\\right]$$\n(17)\n\nThe generalized inverse in formula (17) is represented by []+.\n\n1. (2)\n\nSecondly, matrix B is calculated by formula (18).\n\n$$\\left[\\begin{array}{c}{Y}_{\\mathrm{K}1}\\\\ {}{Y}_{\\mathrm{K}2}\\\\ {}\\mathrm{M}\\\\ {}{Y}_{\\mathrm{K}\\ I}\\end{array}\\right]=\\left[\\begin{array}{c} CdiagA\\left(1,:\\right)\\\\ {} CdiagA\\left(2,:\\right)\\\\ {}M\\\\ {} CdiagA\\left(I,:\\right)\\end{array}\\right]{B}^T+{E}_I$$\n(18)\n\nFormula (18) satisfies the following YK i = Cdiag(A(i,:))BT +EK i, The error is expressed in terms ofEI. The least squares (LS) estimate of BT is calculated by Eq. (19).\n\n$${B}^T={\\left[\\begin{array}{c} CdiagA\\left(1,:\\right)\\\\ {} CdiagA\\left(2,:\\right)\\\\ {}\\mathrm{M}\\\\ {} CdiagA\\left(I,:\\right)\\end{array}\\right]}^{+}\\left[\\begin{array}{c}{Y}_{\\mathrm{K}1}\\\\ {}{Y}_{\\mathrm{K}2}\\\\ {}\\mathrm{M}\\\\ {}{Y}_{\\mathrm{K}I}\\end{array}\\right]$$\n(19)\n1. (3)\n\nThirdly, matrix C is calculated by formula (20).\n\n$$\\left[\\begin{array}{l}{Z}_{\\mathrm{K}1}\\\\ {}{Z}_{\\mathrm{K}2}\\\\ {}\\mathrm{M}\\\\ {}{Z}_{\\mathrm{K}J}\\end{array}\\right]=\\left[\\begin{array}{l} AdiagB\\left(1,:\\right)\\\\ {} AdiagB\\left(2,:\\right)\\\\ {}\\mathrm{M}\\\\ {} AdiagB\\left(J,:\\right)\\end{array}\\right]{C}^T+{E}_J$$\n(20)\n\nWhere ZKj = Adiag(B(j,:))CT + EKj, j = 1, 2, , J. The error is expressed in terms of EJ. The least squares estimate of CT is calculated by Eq. (21).\n\n$${C}^T={\\left[\\begin{array}{c} AdiagB\\left(1,:\\right)\\\\ {} AdiagB\\left(2,:\\right)\\\\ {}\\vdots \\\\ {} AdiagB\\left(J,:\\right)\\end{array}\\right]}^{\\div}\\left[\\begin{array}{c}{Z}_{\\dots 1}\\\\ {}{Z}_{\\dots 2}\\\\ {}\\vdots \\\\ {}{Z}_{\\dots J}\\end{array}\\right]$$\n(21)\n1. (4)\n\nFinally, stop running until the result converges or reaches the set number of iterations after repeating the above steps (1)–(3).\n\nMulti-channel vibration signals are collected in this paper to research the fault state of equipment, and a third-order tensor is constructed through continuous wavelet transform. Figure 3 shows the basic structure of the parallel factor analysis decomposition model for fault diagnosis.\n\nThe Nt, Nd, and Nf of the data matrix $${\\mathrm{S}}_{\\left({\\mathrm{N}}_{\\mathrm{d}}\\times {\\mathrm{N}}_{\\mathrm{f}}\\times {\\mathrm{N}}_{\\mathrm{t}}\\right)}$$ are the number of data points, the number of channels, and the frequency step size, respectively.\n\n$${\\hat{S}}_{dft}=\\sum \\limits_{k=1}^{N_k}{a}_{dk}{b}_{fk}{c}_{ik}+{e}_{dft}$$\n(22)\n\nThe key issue of this parallel factorization model is to obtain the matrices A, B, and C. The adk, bfk, and ctk are their elements, in which component k represents an atom. These spatial signals, spectral signals, and temporal signals for each atom are indicated as the ak = {adk}, bk = {bfk} and ck = {ctk}. The “eijk” is the error, which forms error set E of the third order tensor I × J × K. The uniqueness of the solution of parallel factor decomposition for fault diagnosis is guaranteed through rank(A) + rank(B) + rank(C) ≥ 2Nk + 2. The decomposition of formula (22) is achieved by solving $$\\underset{a_{dk}{b}_{jk}{c}_{ik}}{\\min}\\left\\Vert {\\hat{S}}_{dft}-{\\sum}_{k=1}^{N_k}{a}_{dk}{b}_{fk}{c}_{tk}\\right\\Vert$$. The main advantage of this method is that the spectrum decomposition of time-varying vibration signal is unique and the best model can be obtained under the principle of minimum square deviation.\n\nThe basic steps for implementing the multi-scale parallel factorization model of fault diagnosis are as follows .\n\n1. (1).\n\nAfter the vibration signal is collected, the third order tensor is constructed by continuous wavelet transform.\n\n2. (2).\n\nThe number of factor F is determined by the principle of consistency in MATLAB.\n\n3. (3).\n\nInitialization for load matrix B and C.\n\n4. (4).\n\nAfter initializing and running the matrices B and C, the matrix A is estimated by the least square regression algorithm. A = XZ'(ZZ')−1, Z = (bc).\n\n5. (5).\n\nSimilarly, the matrices B and C are estimated.\n\n6. (6).\n\nContinue from step (3) until the result converges or reaches the set number of iterations.\n\n7. (7).\n\nCorresponding results obtained.\n\n### 2.5 Numerical simulation based on parallel factor analysis\n\nSimulation experiments can investigate the characteristic of the results of input signals with different parameters after parallel factor analysis for fault diagnosis. Therefore, the simulation signals are used to simulate the running state of the centrifugal pump to test the method proposed in this article. The simulation signal is shown in the following formula (23).\n\n$$\\mathrm{y}\\left(\\mathrm{t}\\right)=0.01\\ast \\cos \\left(2\\ast \\mathrm{pi}\\ast 400\\ast \\mathrm{t}-5\\right)\\ast \\exp \\left(-0.5\\ast \\left(\\left(\\left(\\mathrm{t}-0.03\\right)/0.003\\right)\\hat{\\mkern6mu} 2\\right)\\right)$$\n(23)\n\nFigure 4 shows the time-domain diagram of the simulated signal. An impact signal appears in the graph at 0.03 μs, which simulates the signal generated when the system fails. It is shown in Fig. 5 that the time-frequency diagram of the simulation signal is obtained by the continuous wavelet transform. It can be seen from the time-frequency diagram of the simulation signal that the dominant frequency of the signal is 400 Hz and the impact signal in the simulation signal appears in the frequency range of 180–400 Hz. For the time-frequency diagram of the simulated signal in Fig. 5, it can also be seen that the impact signal appears at 0.03 μs.\n\nThe simulation signal is constructed into a third-order tensor after continuous wavelet transformation, and then the third-order tensor is decomposed by parallel factors to obtain the result shown in Fig. 6. After parallel factor analysis, we can get the loading value and residual error corresponding to frequency, time, and channel. Comparing the loading value of frequency and time after parallel factor analysis with the time-frequency diagram, we can find their corresponding relationship. The hypothetical frequency curve in the graph fluctuates in the range of 180–400 Hz, which is a contrast relationship between the fluctuations of the simultaneous frequency graph. The time curve fluctuates at 0.03 μs and has the maximum value of the third component and the minimum value of the second component. It can be seen from the simulation signal corresponding to the ground that the simulation signal also has an impact signal at 0.03 μs. This indicates that the parallel factor analysis for high-dimensional data of multi-source feature factors can well detect the characteristics of the shock signal generated by the simulated fault.\n\nThe vibration signals collected in engineering are generally mixed with various noise signals. In order to check on the effectiveness in complex conditions, we add the noise signal to the original simulation signal and perform parallel factor analysis on it. Figures 7 and 8, respectively, show the time-domain diagram of the original simulation signal after adding noise to the signal and the time-frequency diagram obtained through continuous wavelet transform. After adding the noise signal to the original simulation signal, it can be seen that the waveform of the noisy simulation signal is similar to the original simulation signal in Fig. 4 and the impact signal is almost covered by the noise signal. The waveform of the noisy simulation signal in Fig. 8 is steeper and more rapid, and there is a larger blurred signal at 10–20 Hz.\n\nThe simulation signal with noise is transformed into a third-order tensor after continuous wavelet transformation. The result of the parallel factor analysis of the third-order tensor is shown in Fig. 9. After parallel factor analysis, we can get the loading value and residual parameters corresponding to frequency, time, and channel. Comparing the loading value of frequency and time after parallel factor analysis with the time-frequency diagram, we can find the correspondence between them relationship. The hypothetical frequency curve in the graph fluctuates in the range of 180–400 Hz, which is a contrast relationship between the fluctuations of the simultaneous frequency graph. The time curve fluctuates at 0.03 μs and has a maximum value. We get the normal probability plot and the residual variance corresponding to the data in mode 1, mode 2, and mode 3. This shows that the parallel factor analysis proposed can well detect the characteristics of the impact signal in this paper even when the collected signal contains noise.\n\n## 3 Proposed method\n\nLikelihood function is a function of statistical model parameters, which plays a great role in statistical inference. The general method of using likelihood ratio test statistics was proposed by Neyman-Pearson in 1982 . Its basic idea is similar to the maximum likelihood method of parameter estimation theory, which is called likelihood ratio test. For hypothesis H0 : θ = θ0, alternative hypothesis H1 : θ = θ1, x is a set of random variables. When H0 is true, the probability density function of the random variable x is expressed as f(x, θ0). When H1 is true, the probability density function of the random variable x is expressed as f(x, θ1). The likelihood function of the sample is the following formula (24).\n\n$$L\\left(\\theta \\right)=\\prod \\limits_{i=1}^nf\\left({x}_i,\\theta \\right)$$\n(24)\n\nTherefore, the likelihood ratio test is performed to obtain the statistic L in the following formula (25).\n\n$$L=\\frac{L\\left({\\theta}_1\\right)}{L\\left({\\theta}_0\\right)}=\\frac{\\prod \\limits_{i=1}^nf\\left({x}_i,{\\theta}_1\\right)}{\\prod \\limits_{i=1}^nf\\left({x}_i,{\\theta}_0\\right)}$$\n(25)\n\nIf the likelihood ratio L is larger, the parameter θ is more likely to be θ1; it shows that the result may tend to negate H0. On the contrary, if the ratio is smaller, the parameter θ is more likely to be θ0, which indicates that the result may be inclined to accept H0. For a certain limit k, L is defined as shown in the following formula (26).\n\n$$\\varphi (x)=\\left\\{\\begin{array}{cc}1& l>k\\\\ {}0& l\\le k\\end{array}\\right.$$\n(26)\n\nTest φ(x) is called the likelihood ratio test of the above test problem.\n\nNeyman-Pearson proposes a principle to determine the optimal test method: parameter α satisfies formula (27).\n\n$$\\beta \\left(\\theta \\right)\\le \\alpha \\kern0.5em \\forall \\theta \\in {\\Theta}_0$$\n(27)\n\nIn formula (27), β(θ) is the power function of the test, Θ0 is the parameter space of the null hypothesis H0, and θ is the test parameter. Look for a test that satisfies the above formula so that β(θ) is as large as possible when θ Θ0.To ensure that the probability of making two types of errors is very small, the sample size must be increased. For field testing, the smaller the sample size, the better when ensuring the reliability of the conclusion. The sequential method proposed by A. Wald solves the problem of optimal selection of sample size and play an important milestone in the history of statistical development.\n\nThe probability function f(x, θ) represents the distribution of the random variable x, H0(θ = θ0) and H1(θ = θ1) are the null hypothesis and alternative hypothesis of the random variable x, respectively. When accepting H1, the probability of the sample x1, …, xm for any positive integer m is given by P1m = f(x1, θ1), …, f(xm, θ1), and the probability is given by P0m = f(x1, θ0), …, f(xm, θ0) when accepting H0.The definition of the sequential probability ratio test is as follows: select two normal numbers A and B (B < A) and calculate the probability ratio P1m/P0m at each stage of the test.\n\n1. (a)\n\nIf p1m/p0m ≥ A, the sequential probability ratio test ends, H1is accepted and H0 is discarded.\n\n2. (b)\n\nIf p1m/p0m ≤ B, the sequential probability ratio test ends, H0 is accepted and H1 is discarded.\n\n3. (c)\n\nIf B < p1m/p0m < A, we continue to observe the sequential probability ratio test until the requirement is met.\n\nWhen SPRT is applied to target recognition, it is first assumed that one of the M alternative hypotheses is the initial hypothesis. The signal propagation waveform is denoted as s(t). When a signal is transmitted, one of the possible waveforms is received and recorded as follows:\n\n$$y(t)=s(t)\\ast {h}_i(t)+n(t)\\kern0.5em i\\in \\left\\{1,2,\\mathrm{K},M\\right\\}$$\n(28)\n\nWhere n(t) is additive white Gaussian noise; the impulse response of the target hypothetical channel is expressed as hi(t) and “*” is the convolution factor.\n\nThe signal channel receiving data is defined in formula (29), where Qi represents the target convolution matrix defined in the literature.\n\n$$y={Q}_is+n$$\n(29)\n\nThe M target hypotheses are denoted as H1, H2, …, HM, respectively. The parameter αi,j is the probability (i ≠ j) when the true hypothesis Hi is wrongly selected as Hi. After obtaining k-th observations, suppose the likelihood ratio of i and j can be defined as shown in formula (30)\n\n$${\\Lambda}_{i,j}^k=\\frac{p_{i1}\\left({\\mathrm{y}}_1\\right){p}_{i2}\\left({\\mathrm{y}}_2\\right)\\Lambda \\kern0.5em {p}_{ik}\\left({\\mathrm{y}}_k\\right)\\kern0.5em {P}_i}{p_{j1}\\left({\\mathrm{y}}_1\\right){p}_{j2}\\left({\\mathrm{y}}_2\\right)\\Lambda \\kern0.5em {p}_{j2}\\left({\\mathrm{y}}_k\\right)\\kern0.5em {P}_j}$$\n(30)\n\nWhere pik(yk) is the probability density function (PDF) with k-th data under the i-th hypothesis and yk is the k-th observation data. When the likelihood ratio satisfies formula (31), accept the assumption Hm.\n\n$${\\Delta}_{i,j}^k>\\frac{1-{\\theta}_{i,j}}{\\theta_{i,j}}\\kern0.5em j\\ne i$$\n(31)\n\nWhen the likelihood ratio satisfies the formula (31), stop the loop. If the likelihood ratio does not meet the stopping condition, continue to the next iteration. In fact, the probability density function of the observed data is constant and satisfies pi1(y) = pi2(y) = … = pik(y). The intensity waveform is updated with the number of iterations, so the probability density function of the observation data under the condition of additive white Gaussian noise can be defined as formula (32).\n\n$${p}_{ik}\\left({y}_k\\right)=\\frac{1}{{\\left(\\sqrt{2{\\pi \\sigma}_n^2}\\right)}^{L_y}}\\times \\exp \\left[-\\frac{1}{2{\\sigma}_n^2}{\\left({y}_k-{Q}_i{s}_k\\right)}^T\\kern0.5em \\left({y}_k-{Q}_i{s}_k\\right)\\right]$$\n(32)\n\n## 4 Experiments\n\n### 4.1 Slurry pump fault test system and experimental design\n\nThe experimental system to be established in this project is required to operate the slurry pump under controlled conditions of speed, flow rate, slurry density, and inlet pressure, and to use and replace the impeller of the slurry pump of different grades and wear. Common failure parts of centrifugal pumps include rotor impeller, rolling bearing, seal, coupling, etc., of which impeller and rolling bearing failure account for a large proportion. The schematic diagram of the slurry pump fault diagnosis test system is shown in Fig. 10. The figure shows the three-dimensional schematic diagram of the test circuit and identifies the key components. It mainly includes motors, centrifugal pumps, data acquisition systems, control instruments, glycol cooling tanks, pressure gauges, flow meters, conveyor belts, sand tanks, pipelines, pressure control tanks, density meters, and sampling ports. First, the normal impeller is used in the centrifugal pump to run the slurry pump fault test system for collecting and testing the signal data of the slurry pump vibration, flow, slurry density, motor speed, and pump inlet and outlet pressure. Then impeller perforation, impeller edge damage and blade damage, and its impellers with different degrees of damage were selected to replace the original centrifugal pump impeller. After running the slurry pump fault diagnosis and test system, the data of the vibration, speed, and pump speed of the slurry pump experiment system were collected.\n\nFigure 11 shows the process flow chart of the slurry pump fault diagnosis test system. The arrow direction in the figure is the flow direction of the mud when the mud pump fault diagnosis experiment system is running. It is the basis for establishing and running the centrifugal pump fault diagnosis experiment system in this article. The serial number and related schematic diagram in Figure 11 indicate the following meanings: 1—centrifugal pump, 2—motor, 3—inverter, 4—power meter sensor, 5—accelerometer, 6—pressure sensor, 7—flow meter, 8—hole plate, 9—heat exchanger, 10—cooler, 11—temperature sensor, 12—sand, 13—suction pressure control tank, and 14—suction pressure sensor. The fault diagnosis test system for slurry pump contains a Weir/Warman 3/2 CAH slurry pump (40 HP) with impeller C2147(8.4\"). The process flow chart of fault diagnosis test system for slurry pump covers the key issues mentioned in this article, but does not cover all aspects of the design of the experimental system. The key issues include that the medium of the cooler in the pipeline is ethylene glycol, the process water is municipal water, and the heat exchanger medium is steam. Microphone means for sound collector.\n\nIn order for the experiment to run successfully, the designer first needs to design the system after engineering calculation and determine the components. The main equipment required for the experiment includes the centrifugal pump, data acquisition system for vibration data acquisition, sensors, and a laptop computer. Auxiliary equipment including storage tanks, valves, instruments, and drive motors are used to control various functions. The data collected by vibration accelerometers is used to analyze the centrifugal pump system in this experiment. The detailed explanation of the vibration sensor for the system signal acquisition is shown below. In the experiment, three three-axis vibration accelerometers are used. Two of the PCB three-axis ICP (Integrated Circuit Piezoelectric) accelerometers have the sensitivity of 100 mV/g and the frequency range of 2–5 kHz. Another PCB three-axis ICP accelerometer has the range of 0.5–3 kHz and the sensitivity of 1000 mV/g.\n\n### 4.2 Slurry pump experimental equipment and signal acquisition system\n\nTo research the validity of the multi-scale parallel factor analysis and sequential probability ratio test proposed in this paper for multi-channel signal in actual complex industrial production, the centrifugal pump fault diagnosis experimental system was designed. The general Fig. 12a shows the centrifugal pump fault diagnosis experimental system. The data acquisition system is shown in Fig. 12b based on a combination of PC measurement hardware and software, which can input electrical signals from sensors and other instruments into a computer for processing. NI LabView 7.0 was chosen as the measurement standard application software because it is easy to build a graphical measurement interface with the help of a large number of tools and objects. The selected hardware is provided by NI DAQ and is highly compatible with our software applications. In order to collect the vibration signals of the centrifugal pump in three directions for each state, it is necessary to install a short-range but high-sensitivity sensor at the key position. Figure 12d is a schematic diagram of the position of the accelerometer. The standard accelerometer and the high-sensitivity accelerometer are installed on the pump casing near the pump suction port, where they will be close to the parts that are prone to failure. Another standard accelerometer is mounted on the shaft bearing because this location is sensitive to vibrations transmitted from the stuffing box. Real-time signals such as flow rate, pressure, speed, and vibration, can be collected synchronously by the experimental data acquisition system. By commanding the pressure and the flow of the equipment’s loop, we simulate the non-linear operating state of the industrial process of the mechanical system to establish the non-linear multi-fault mode, synchronously collect multi-channel signals, and obtain multi-source signals. The internal interaction mechanism between fluid excitation and vibration response under non-linear operation mechanism can be analyzed.\n\nThe liquid transported in this experiment is set as mud to better collect vibration signals. In this experiment, normal impeller, and three types of faulty impellers, including impeller perforation, impeller edge damage and blade damage, were set to simulate failures in industrial production. Among them, these three failure modes have clear differences and the typical failures of centrifugal pump impellers can be represented well. The impeller in the normal state is denoted as S1, and the three types of impellers with impeller perforation, impeller edge damage, and blade damage are denoted as S2, S3, and S4. In order to avoid aliasing, the sampling frequency in this experiment is 9009 Hz according to the Nyquist sampling theorem and the data acquisition time is 20 s for each group.\n\nIn the experiment, different impellers were replaced to collect the vibration signals of the centrifugal pump under different operation conditions. The steps of the whole experiment are summarized as follows:\n\n1. (1)\n\nEstablish the experimental device according to the schematic diagram of slurry pump test system shown in Fig. 11. The normal impeller shown in Fig. 13 is used as the impeller of the pump and sediment is added as the pumping medium. After starting the motor, we adjust the motor speed to 1200 rpm, 1600 rpm, 1800 rpm, 2200 rpm, and 2600 rpm through the known voltage, motor power, motor efficiency, and other coefficients and instructions. According to the sampling time of 20 s and sampling frequency of 9009 Hz shown in the previous article, NI LabView 7.0 application software and NI DAQ signal acquisition system were run to collect the three sets of three-dimensional vibration signals of the corresponding pump.\n\n2. (2)\n\nThe normal impeller in the original centrifugal pump is replaced by the impeller perforation of the fault S2 in Fig. 13, and the other parts remain unchanged. Follow the previous steps to start the centrifugal pump and collect data. When a set of data is collected, the speed is set to 1400 rpm, 1600 rpm, and 2600 rpm and the above steps are repeated to collect data.\n\n3. (3)\n\nThe S3 of impeller edge damage in Fig. 13c is selected to replace the impeller of S2 in the original centrifugal pump and other parts remain unchanged. Similarly, follow the previous steps to start the centrifugal pump and collect data.\n\n4. (4)\n\nThe S4 of blade damage in Fig. 13d is selected to replace the impeller of S3 in the original centrifugal pump and other parts remain unchanged. Similarly, follow the previous steps to start the centrifugal pump and collect data.\n\n5. (5)\n\nAfter the experiment, the outlet valve of the pump was closed. Close the inlet valve after turning off the motor. Store experimental data to prepare for subsequent vibration signal analysis.\n\n## 5 Results and discussion\n\n### 5.1 Multi-source dynamic feature extraction based on parallel factorization\n\nThis multi-scale parallel factorization method for the extraction of characteristic signals in non-linear multi-source and multi-fault modes is proposed in the article. Parallel factorization can not only perform high-dimensional data processing, but also has the uniqueness of the decomposition. This property makes the results of parallel factorization more realistic and has specific physical meanings. The third-order tensor constructed by multi-channel vibration signals through continuous wavelet transform is decomposed into the series of different modes of channel/frequency/time by the multi-scale parallel factor analysis algorithm. The spatial information is introduced into the time-frequency analysis of signals to form the three-dimensional spatial/time/frequency characteristic analysis of each factor. The simulation results show that the parallel factor decomposition for the tensor built by multi-channel signal has the compatibility of decomposition path and overall consistency. As a result, the topographic map, spectrum, and time contour of the multi-source fault signal in the centrifugal pump experiment are acquired. The multi-scale parallel decomposition method for extracting multi-source feature signals of non-linear failure modes is applied in the fault diagnosis of centrifugal pumps. It analyzes the internal connection between the optimal decomposition paths of multi-source signal feature factors. The optimal non-linear correspondence relationship between failure modes and characteristic signals in time, frequency, and space are constructed. Based on the correspondence and overall consistency of multi-source feature factor decomposition paths, we remodeled three-dimensional fault feature models such as the frequency spectrum and time profile of the fault feature factors, successfully extracted non-linear multi-dimensional dynamic fault feature signals. Finally, the corresponding fluctuation regularities of the homologous non-linear failure mode in the multi-source signals were displayed.\n\nFigure 14 shows the time-frequency diagram obtained of the vibration signals collected in the X-axis direction of the three vibration signal collection points by continuous wavelet transformation when the slurry pump is in normal operation. Figure 15 shows the result of the parallel factor analysis of vibration signal of slurry pump after continuous wavelet transform in normal state. In this experiment, three-dimensional vibration sensors are set up at three measuring points. We analyze the vibration signals of these three measuring points to explore the three-dimensional spatial distribution and characteristic propagation path of dynamic characteristics on the mechanical structure of the slurry pump. Three groups of original vibration signals are transformed by continuous wavelet to obtain three-dimensional time-frequency signals to construct a third-order matrix. After multi-scale parallel factor analysis for the third-order tensor, the loading values and residual variance of the aisle, time, and frequency factors are obtained.\n\nFigure 16 shows the time-frequency diagram of the vibration signals collected in the X-axis direction of the three vibration signal collection points by continuous wavelet transformation when the slurry pump is in S2 impeller perforation. In the S2 state, the third-order tensor of 3 × 126 × 4096 is constructed by continuous wavelet transform. Figure 17 indicates the result of the loading values and residual variance of the aisle, time, and frequency modes by the parallel factor analysis for the third-order tensor of slurry pump in state S2.\n\nFigure 18 shows the time-frequency diagram of the vibration signals collected in the X-axis direction of the three vibration signal collection points by continuous wavelet transformation when the slurry pump is in S3 impeller edge damage. In the S3 state, the third-order tensor of 3 × 126 × 4096 is constructed by continuous wavelet transform. Figure 19 indicates the result of the loading values and residual variance of the aisle, time, and frequency modes by the parallel factor analysis for the third-order tensor of slurry pump in state S3.\n\nThe operating state of the slurry pump system in our experimental device system has normal operation and three failure states. The failure states include impeller holes, leading edge damage, and propeller blade damage. Similarly, Fig. 20 shows the time-frequency diagram when the slurry pump is in S4 propeller blade damage. Figure 21 indicates the result of the loading values and residual variance of the aisle, time and frequency modes by the parallel factor analysis for state S4.\n\nWe analyze the vibration signals of these three measuring points to discuss the three-dimensional spatial distribution and characteristic propagation path of dynamic characteristics on the mechanical structure of the slurry pump. By comparing the decomposition results of parallel factor analysis in the normal and fault state, there are obvious difference in the time loading factor and frequency loading factor component. Due to the phenomenon of characteristic coupling and aliasing of mechanical multi-source signals, the parallel factor analysis can optimize the independent characteristics of each channel on the surface of the mechanical structure and eliminate the mutual interference, overlap, and redundancy of the characteristic signals between the channels. Therefore, the parallel factor analysis are effective in providing a basis for subsequent diagnosis of SPRT and the fault identification can be successfully implemented.\n\n### 5.2 SPRT for the multi-source condition monitoring of centrifugal pump\n\nThe proportions of the standard deviation “σ” and mean “μ” of the test signal sequences have significant influence on the likelihood ratio in the sequential probabilistic ratio test. Therefore, the mean value and standard deviation of the frequency loading value after the parallel factor decomposition should be calculated first for the test signal sequence. Assuming the probability distribution of the frequency load value sequence of one set of signals under the multi-source condition monitoring of the centrifugal pump meets the null hypothesis Hi : μ = μi, and the probability distribution of the frequency load value sequence of the other set of signals satisfies the alternative hypothesis Hj : μ = μj . Their corresponding standard deviation σ remains unchanged. When the original hypothesis and the alternative hypothesis are both true, the joint probability density functions of these two sets of sequences are shown below.\n\n$${p}_{ik}\\left({y}_k\\right)=\\frac{1}{\\sigma \\sqrt{2\\pi }}\\exp \\left(-\\frac{1}{2{\\sigma}^2}{\\left({y}_k-{\\mu}_i\\right)}^2\\right)$$\n(33)\n$${p}_{jk}\\left({y}_k\\right)=\\frac{1}{\\sigma \\sqrt{2\\pi }}\\exp \\left(-\\frac{1}{2{\\sigma}^2}{\\left({y}_k-{\\mu}_i\\right)}^2\\right)$$\n(34)\n\nIn formula (33), Pik(yk) is the probability density function null hypothesis. Pjk(yk) in formula (34) is the probability density function under the alternative hypothesis. The SPRT probability ratio is calculated in formula (35).\n\n$${\\lambda}_{i,j}\\left({Y}_{Sm}\\right)=\\frac{\\prod \\limits_{k=1}^n{P}_{jk}}{\\prod \\limits_{k=1}^n{P}_{ik}}=\\frac{P_{j1}\\left({y}_1\\right){P}_{j2}\\left({y}_2\\right)\\Lambda \\kern0.5em {P}_{jk}\\left({y}_k\\right)}{P_{i1}\\left({y}_1\\right){P}_{i2}\\left({y}_2\\right)\\Lambda \\kern0.5em {P}_{ik}\\left({y}_k\\right)}\\times \\frac{P_{j0}}{P_{i0}}$$\n(35)\n\nIn order to make the calculation easier in practical applications, the likelihood ratio formula is further derived and simplified to obtain the formula (36). Where YSi and YSj are the to-be-checked sequences of vibration signals Si and Sj, respectively, Δi, j(YSi) and Δi, j(YSj) are the likelihood ratios of the sequence to be tested YSi and YSj, respectively.\n\n$${\\Delta}_{i,j}\\left({Y}_{Sm}\\right)=1\\mathrm{n}{\\lambda}_{i,j}\\left({Y}_{Sm}\\right)=1\\mathrm{n}\\frac{\\prod \\limits_{k=1}^n{P}_{jk}}{\\prod \\limits_{k=1}^n{P}_{ik}}=\\sum \\limits_{k=1}^n1\\mathrm{n}\\frac{P_{jk}}{P_{ik}}\\kern0.5em m=i,j$$\n(36)\n\nReferring to the sequential probability ratio test algorithm, we compare the likelihood ratio with the thresholds A and B to identify different forms of failure of the centrifugal pump. The size of A and B are closely related to the probability α of type I error and the probability β of type II error. The variables α, β, A, and B are satisfied with the following relationship:\n\n$$a=\\ln A=\\ln \\frac{1-\\beta }{\\alpha }$$\n(37)\n$$b=\\ln B=\\ln \\frac{\\beta }{1-\\alpha }$$\n(38)\n\nFor S1, S2, S3, and S4 four different impeller fault states of the centrifugal pump experimental system, Figure 22 shows the process of using the sequential probability ratio test algorithm to identify the fault. Vibration signals at three different positions collected under two different impeller fault conditions (sj) and (sj) are decomposed by parallel factors to obtain frequency loading values, which are calculated according to formulas (33)–(36) to obtain the likelihood ratio Δi, j of the sequential probability ratio test. The process of centrifugal pump fault identification is shown below. (1) If Δi, j =  ln (λi, j)  (−∞, b], accept Hj, the centrifugal pump system is under the condition (sj). (2) If Δi, j =  ln (λi, j)  [a, ∞), accept Hi, the centrifugal pump is under the condition (si). (3) If Δi, j [a,  b], the likelihood ratio of sequential probabilistic ratio test continues to be calculated by extracting the next data in the test sequence according to formulas (33)–(36). The likelihood ratio will continue to be compared with the threshold value until the condition (1) or (2) is met or the number of iterations is reached. After the test is stopped and the probability parameters λ1, 2(YS1), λ1, 2(YS2), λ1, 3(YS1), λ1, 3(YS3), λ1, 4(YS1) andλ1, 4(YS4) are obtained, the conditions of centrifugal pump will be distinguished.\n\nThe means of the signal S1, S2, S3 and S4 under the four conditions parameters areμ1, μ2, μ3, μ4. Then, the likelihood ratio is calculated and analyzed according to formulas (33)~(36). SPRT probability ratios λi, j(YSi) and λi, j(YSj) are calculated by importing the testing data (YSi, YSj) of the signal waveform for slurry pump Si and Sj conditions to Eq. (35). Compare the likelihood ratios Δi, j(YSi) and Δi, j(YSj) with the threshold to determine the state Si and Sj of the centrifugal pump.\n\nThe difference in the variables λi, j(YSi) and λi, j(YSj) is used to distinguish different condition (Si,Sj) of the centrifugal pump. When the iteration periods are determined, the condition S1 is compared with S2, S3, S4 in the relation between SPRT probability ratios. Figure 23a shows the fluctuation curve between the likelihood ratios Δ1, 2(YS1) and Δ1, 2(YS2) of the signals S1 and S2 during the determined number of iteration. It can be seen from the curve in Fig. 23a that Δ1, 2(YS1) < b displays that the centrifugal pump is in the normal state of S1; Δ1, 2(YS2) > a indicates that the centrifugal pump is in a fault state of the S2 impeller perforation. From the curve in Fig. 23b, it can be seen that Δ1, 3(YS1) < b indexes that the centrifugal pump is in the normal state of S1; Δ1, 3(YS3) > a indicates that the centrifugal pump is in the fault state of S3 impeller edge damage. In Fig. 23c, the inequality (Δ1, 4(YS1)) < b is satisfied, the condition is S1. The inequality(Δ1, 4(YS4)) > a is satisfied, the condition of the centrifugal pump is S4 impeller blade damage. The SPRT parameters Δ1, m, m = 2, 3, 4 in Fig. 22a–c are adopted to distinguish that the normal condition (S1) of the centrifugal pump from the conditions (S2,S3,S4).\n\nWe found that the different fault states of the centrifugal pump in the experiment can also be distinguished by this method of SPRT. Figure 24a indicates the SPRT parameters λ2, 3 in Eqs. (33–36) of the testing sequences YS2 and YS3. The SPRT inequality (Δ2, 3(YS2)) < b is satisfied, then the condition of the centrifugal pump is S2. The inequality (Δ2, 3(YS3)) > a is satisfied, then the condition of the slurry pump is S3 (impeller edge damage). It can be seen from the curve in Fig. 24b that when (Δ2, 4(YS2)) < b is satisfied, the centrifugal pump is fault S2 (impeller perforation); when (Δ2, 4(YS4)) > a is satisfied, the centrifugal pump is fault S4 (impeller blade damage). When the inequality (Δ3, 4(YS3)) < b is satisfied in Fig. 24c, the condition is S3 (impeller edge damage). When (Δ3, 4(YS4)) > a is satisfied, the condition of the centrifugal pump is S4 (impeller blade damage). The parameters (λi, j(YSi), λi, j(YSj)) are effective indicator to monitor the different conditions of the multi-fault and multi-source centrifugal pump.\n\n## 6 Conclusion\n\nParallel factorization can not only perform high-dimensional data processing, but also has the uniqueness of the decomposition. This property makes the results of parallel factorization more realistic and has specific physical meanings. Through numerical simulation, the parallel factorization is used to explore the non-linear correspondence relationship of multi-source signal characteristic factors in time, frequency, and space under different simulation states. By adjusting the different frequency, time, phase, and amplitude of the analog signal, the loading values of the three modes are captured after parallel factorization to research the corresponding relationship between the analog signals. Then, the parallel factor analysis is applied to the centrifugal pump fault diagnosis experimental system to analyze the state characteristics under multiple fault modes. The non-linear multi-dimensional actional fault characteristic parameter of the impellers with different faults of the centrifugal pump was triumphantly acquired, and the corresponding fluctuation regularities of the homologous fault mode characteristics in the multi-source signals were displayed.\n\nThe analysis of the comprehensive result graph shows that the centrifugal pump fault diagnosis methodology based on parallel factor analysis of multiple scales and sequential probability ratio test is effective and reliable. This method first analyzes the collected vibration signals through parallel factor analysis and then conducts sequential probability ratio testing. It identifies different failure modes by comparing likelihood ratios and thresholds. Not only normal conditions and fault conditions can be identified, but also different fault conditions can be distinguished. Therefore, the methodology proposed in these contents of article is very suitable for non-linear multi-source and multi-fault signal analysis and processing. The PARAFAC theory proposed in this paper can also be used in the blind separation of mechanical multiple faults.\n\n## Availability of data and materials\n\nAll data generated or analyzed during this study are included in this published article.\n\n## Abbreviations\n\nPARAFAC:\n\nParallel factor analysis\n\nSPRT:\n\nSequential probability ratio test\n\nCNN:\n\nConvolutional neural network\n\nEMD:\n\nEmpirical mode decomposition\n\nVMD:\n\nVariational mode decomposition\n\nEEG:\n\nElectroencephalogram\n\nANN-MLP:\n\nArtificial neural network-multi-layer perceptron\n\nTALS:\n\nTrilinear alternate least squares\n\nLS:\n\nLeast squares\n\nPDF:\n\nProbability density function\n\n## References\n\n1. R. Azizi, B. Attaran, A. Hajnayeb, A. Ghanbarzadeh, M. 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Boris, et al., Application of Parallel Factor Analysis (PARAFAC) to electrophysiological data. Front. Neuroinformatics 8 (2015). https://doi.org/10.3389/fninf.2014.00084.\n\n25. F. Miwakeichi, E. Martinez-Montes, et al., Decomposing EEG data into space-time-frequency components using Parallel Factor Analysis. Neuroimage 22(3), 1035–1045 (2004)\n\n26. Z. Rost‘akova, R. Rosipa, S. Seifpour, et al., A Comparison of Non-negative Tucker Decomposition and Parallel Factor Analysis for Identification and Measurement of Human EEG Rhythms. Meas. Sci. Rev. 20(3), 126–138 (2020)\n\n27. C.J. Yeh, H. Heungsun, E. Timmerman Marieke, Functional Parallel Factor Analysis for Functions of One- and Two-dimensional Arguments. Psychometrika 83(1), 1–20 (2018)\n\n28. W. Jing, W. Juan, W. You, L. Xiaoli, The frontal area with higher frequency response is the principal feature of laser-evoked potentials in rats with chronic inflammatory Pain: a Parallel Factor analysis study. Front. Neurol. 8 (2017). https://doi.org/10.3389/fneur.2017.00155.\n\n29. L. Yang, H. Chen, Y. Ke, L. Huang, Q. Wang, Y. Miao, L. Zeng, A novel time-frequency-space method with parallel factor theory for big data analysis in condition monitoring of complex system. Int. J. Adv. Robotic Syst. 17(2) (2020). https://doi.org/10.1177/1729881420916948.\n\n30. Y. Cheng, J. Minping, A novel weak fault signal detection approach for a rolling bearing using variational mode decomposition and phase space parallel factor analysis. Meas. Sci. Technol. 30(11) (2019). https://doi.org/10.1088/1361-6501/ab30bd.\n\n31. L. Lijia, B. Shiyi, M. Jianfeng, et al., Monitoring batch processes using sparse parallel factor decomposition. Ind. Eng. Chem. Res 56(44), 12682–12692 (2017)\n\n32. G. Haiyang, L. Kaiqi, H. Xingyi, et al., Feasibility study for the analysis of coconut water using fluorescence spectroscopy coupled with PARAFAC and SVM methods. Br. Food J (2020). https://doi.org/10.1108/BFJ-12-2019-0941\n\n33. H. Xiu, Y. Huibin, S. Yonghui, et al., Characterizing humic substances from a large-scale lake with irrigation return flows using 3DEEM-PARAFAC with CART and 2D-COS. J. Soils Sediments 20(9), 3514–3523 (2020)\n\n34. C. Yahya, A sequential probability ratio test (SPRT) to detect changes and process safety monitoring. Process Saf. Environ. Protect. 92(3), 206–214 (2014)\n\n35. G. Peng, I. David, Wind turbine power curve modeling and monitoring with Gaussian process and SPRT. IEEE Transact. Sustainable Energy 11(1), 107–115 (2020)\n\n36. W. Rong, X. Zhi, L. Jianye, et al., Chi-square and SPRT combined fault detection for multisensor navigation. IEEE Transact. Aerospace Electron. Syst. 52(3), 1352–1365 (2016)\n\n37. H. Kasai, Fast online low-rank tensor subspace tracking by CP decomposition using recursive least squares from incomplete observations. Neurocomputing 347, 177–190 (2019)\n\n38. W. Huang, J. Huang, L. Yang, et al., Fault diagnosis of gearboxbased on principal component analysis and sequential probability ratio test [C]. Proceedings of 2nd International Conference on Computer Science and Application Engineering, Hohhot, China. (2018). https://doi.org/10.1145/3207677.3277945\n\n39. P. Juan, Z. Santiago, Using one class SVM to counter intelligent attacks against an SPRT defense mechanism. Ad Hoc Networks. 94 (2019). https://doi.org/10.1016/j.adhoc.2019.101946.\n\n## Acknowledgements\n\nThe Reliability Research Lab in the Department of Mechanical Engineering at the University of Alberta in Canada provided the original experimental data. The National Natural Science Foundation of China (Grant No. 51775390, 51901164, 51805378), the Natural Science Foundation of Hubei Province (Grant No. 2018CFB394), and the Foundation of Wuhan Science and Technology Bureau (Grant No. 2019010701011417) provides the financial support for paper research.\n\n## Funding\n\nThis research was funded by the National Natural Science Foundation of China (Grant No. 51775390, 51901164, 51805378), the Natural Science Foundation of Hubei Province (Grant No. 2018CFB394), and the Foundation of Wuhan Science and Technology Bureau (Grant No. 2019010701011417).\n\n## Author information\n\nAuthors\n\n### Contributions\n\nThe algorithms proposed in this paper have been conceived by all authors. Hanxin Chen designed and performed the experiments and then analyzed the results. Liu Yang performed the experiments and analyzed the simulation results. Liu Yang wrote the paper. All authors read and approved the final manuscript.\n\n### Corresponding author\n\nCorrespondence to Hanxin Chen.\n\n## Ethics declarations\n\nNot applicable.\n\nNot applicable.\n\n### Competing interests\n\nThe authors declare that they have no competing interests.", null, "" ]
[ null, "https://asp-eurasipjournals.springeropen.com/track/article/10.1186/s13634-021-00730-w", null ]
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https://convertoctopus.com/39-feet-per-second-to-knots
[ "## Conversion formula\n\nThe conversion factor from feet per second to knots is 0.59248380129641, which means that 1 foot per second is equal to 0.59248380129641 knots:\n\n1 ft/s = 0.59248380129641 kt\n\nTo convert 39 feet per second into knots we have to multiply 39 by the conversion factor in order to get the velocity amount from feet per second to knots. We can also form a simple proportion to calculate the result:\n\n1 ft/s → 0.59248380129641 kt\n\n39 ft/s → V(kt)\n\nSolve the above proportion to obtain the velocity V in knots:\n\nV(kt) = 39 ft/s × 0.59248380129641 kt\n\nV(kt) = 23.10686825056 kt\n\nThe final result is:\n\n39 ft/s → 23.10686825056 kt\n\nWe conclude that 39 feet per second is equivalent to 23.10686825056 knots:\n\n39 feet per second = 23.10686825056 knots\n\n## Alternative conversion\n\nWe can also convert by utilizing the inverse value of the conversion factor. In this case 1 knot is equal to 0.04327717582307 × 39 feet per second.\n\nAnother way is saying that 39 feet per second is equal to 1 ÷ 0.04327717582307 knots.\n\n## Approximate result\n\nFor practical purposes we can round our final result to an approximate numerical value. We can say that thirty-nine feet per second is approximately twenty-three point one zero seven knots:\n\n39 ft/s ≅ 23.107 kt\n\nAn alternative is also that one knot is approximately zero point zero four three times thirty-nine feet per second.\n\n## Conversion table\n\n### feet per second to knots chart\n\nFor quick reference purposes, below is the conversion table you can use to convert from feet per second to knots\n\nfeet per second (ft/s) knots (kt)\n40 feet per second 23.699 knots\n41 feet per second 24.292 knots\n42 feet per second 24.884 knots\n43 feet per second 25.477 knots\n44 feet per second 26.069 knots\n45 feet per second 26.662 knots\n46 feet per second 27.254 knots\n47 feet per second 27.847 knots\n48 feet per second 28.439 knots\n49 feet per second 29.032 knots" ]
[ null ]
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https://electronics.stackexchange.com/questions/415955/how-does-this-bjt-affect-the-circuit-if-it-does-how-should-i-simulate-on-circu?noredirect=1
[ "# How does this BJT affect the circuit? If it does, how should I simulate on Circuitlab to see the effect?\n\nI am trying to understand how NPN Q23 in the picture below affects the operation of the left circuit.", null, "If it doesn't then what we get is the circuit on the right with the only transistor now being Q30.\n\nAfter reading this thread I would say Q23 will be in reverse active mode of operation operating with lower β and therefore it does affect the operation of the circuit on the left.\n\nMy questions are:\n\n1) What is it that actually happens by placing Q23 on the circuit? Is there more current passing through R1 because the resistance of Q23//Q22 is smaller?\n\n2) If I am to simulate the circuit (on Circuitlab or Multisim) what simulation setting (DC sweep, time domain) and input parameters would you suggest so that I can see a clear difference between the left and the right circuit?\n\nEDIT:\n\nThank you for the comments. This is a made up question, I am interested in the effect Q23 will have on the ADC voltage across R1 given that V2 is constant DC source. I simply want to establish that Q23 and Q22 will have the same constant base current and see what happens from there..\n\nEDIT#2\n\nThe question could very well be about the effect of Q23 being the only BJT in the left circuit. In that case what simulation would you run to compare the circuits on the left and right?\n\n• What do voltage sources V1 and V2 represent? Are they AC, DC, pulses, sensor outputs, what? – AnalogKid Jan 8 at 22:57\n• This might be a good question if you'd provide more context about it. (If it is just made up, then that's another thing.) Because of the heavier doping of the emitter, there is an avalanche behavior in $Q_{23}$ that might be sought. $Q_{23}$'s crappy $\\beta$ means $R_2$ will have to drop more voltage. $Q_{23}$'s forward-biased voltage drop will be lower and pinch $Q_{22}$ still more, so the aggregate $\\beta$ might be something desired. But what to suggest without any context? No idea. – jonk Jan 9 at 3:43\n• @jonk and thank you for the answers. I edited the question adding some further context, hope it clarifies a few things. – George Jan 9 at 14:19\n• @AnalogKid V1, V2 are constant DC sources. Please read the edit for clarification. Thanks for the comment! – George Jan 9 at 14:20\n• I know that this is what you want to measure. But sometimes we add components for protection or sideeffects. – le_top Jan 10 at 22:41\n\nYou are right the Q23 will be ON. Why? Because there is a path for a base-collector current to flow for NPN transistor to ground.", null, "simulate this circuit – Schematic created using CircuitLab\n\nAnd in this reverse connector the collector took over the role of the emitter and the emitter now takes the role of the collector. But due to the different doping (and size) between the collector and the emitter. The reverse beta $$\\\\beta_R\\$$ will be much lower than the \"normal\" forward beta $$\\\\beta_F\\$$.\n\nFor example my BC548B show this resoult: $$\\\\beta_F = 250\\$$ at $$\\1\\textrm{mA}\\$$ and $$\\\\beta_R = 8.3\\$$ in reverse active mode for the same current.\n\nAnd for BC337-25 $$\\\\beta_F = 352\\$$ ; $$\\\\beta_R = 38\\$$\n\nAs a side note the BJT will also conduct current in these two cases:\n\nBJT as a Zener diode", null, "BJT's now behaviors just like a poor's man tunnel diode (Esaki effect). And tunnel diodes will have \"negative resistor\" region. And this negative resistance region occurs only for NPN BJT's.", null, "http://jlnlabs.online.fr/cnr/negosc.htm\n\nhttp://www.cappels.org/dproj/simplest_LED_flasher/Simplest_LED_Flasher_Circuit.html\n\nBC337-40\n\nVeb=8.2V , Vec=6.7V at I=5.5mA\n\nBC549B\n\nVeb=8.3V , Vec=7.2V at I=5.5mA\n\nBD139-16\n\nVeb=8.5V, Vec=6.7V at I=5.5mA\n\n• But how is the base-emitter (PN) diode forward biased? Base is at 1 V and emitter at 3 V. – thece Feb 1 at 18:09\n• @thece But in this circuit, the BJT is in reverse-active mode. Which means that the Base-Collector is in forward biased and Base-Emitter is in reverse biased. – G36 Feb 1 at 19:16" ]
[ null, "https://i.stack.imgur.com/C7qXb.png", null, "https://i.stack.imgur.com/smnfQ.png", null, "https://i.stack.imgur.com/AtAxX.png", null, "https://i.stack.imgur.com/vg54l.png", null ]
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https://www.ams.org/journals/tran/1973-180-00/S0002-9947-1973-0318115-3/home.html
[ "# Transactions of the American Mathematical Society\n\nPublished by the American Mathematical Society since 1900, Transactions of the American Mathematical Society is devoted to longer research articles in all areas of pure and applied mathematics.\n\nISSN 1088-6850 (online) ISSN 0002-9947 (print)\n\nThe 2020 MCQ for Transactions of the American Mathematical Society is 1.48.\n\nWhat is MCQ? The Mathematical Citation Quotient (MCQ) measures journal impact by looking at citations over a five-year period. Subscribers to MathSciNet may click through for more detailed information.\n\nby R. L. Davis\nTrans. Amer. Math. Soc. 180 (1973), 47-52 Request permission\n\n## Abstract:\n\nLet \\$K\\$ be a field having prime characteristic \\$p\\$. The following conditions on a subfield \\$k\\$ of \\$K\\$ are equivalent: (i) \\${ \\cap _n}{K^{{p^n}}}(k) = k\\$ and \\$K/k\\$ is separable. (ii) \\$k\\$ is the field of constants of an infinite higher derivation defined in \\$K\\$. (iii) \\$k\\$ is the field of constants of a set of infinite higher derivations defined in \\$K\\$. If \\$K/k\\$ is separably generated and \\$k\\$ is algebraically closed in \\$K\\$, then \\$k\\$ is the field of constants of an infinite higher derivation in \\$K\\$. If \\$K/k\\$ is finitely generated then \\$k\\$ is the field of constants of an infinite higher derivation in \\$K\\$ if and only if \\$K/k\\$ is regular.\nSimilar Articles\n• Retrieve articles in Transactions of the American Mathematical Society with MSC: 12F15\n• Retrieve articles in all journals with MSC: 12F15" ]
[ null ]
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https://www.maplesoft.com/support/help/maple/view.aspx?path=Maplets/Elements/RunDialog&L=E
[ "", null, "Maplets[Elements] - Maple Programming Help\n\nHome : Support : Online Help : Programming : Maplets : Elements : Command Elements : Maplets/Elements/RunDialog\n\nMaplets[Elements]\n\n RunDialog\n display a dialog\n\n Calling Sequence RunDialog(opts)\n\nParameters\n\n opts - equation(s) of the form dialog=value; specify options for the RunDialog element\n\nDescription\n\n • The RunDialog command element displays a dialog defined in the running Maplet application. To display a dialog defined in a different Maplet application, use the Display function.\n • There are separate RunDialog and RunWindow command elements because the Window element is intrinsically different from the dialog elements. A dialog has a predefined structure. A Maplet application author can specify options for a dialog, but cannot add elements. A window does not have a predefined structure. A Maplet application author specifies its structure by using elements and options. Also they behave differently. For example, a window can be minimized.\n • To simplify specifying options in the Maplets package, certain options and contents can be set without using an equation. The following table lists elements, symbols, and types (in the left column) and the corresponding option or content (in the right column) to which inputs of this type are, by default, assigned.\n\n Elements, Symbols, or Types Assumed Option or Content name or string dialog option\n\n • A RunDialog element cannot contain other elements.\n • A RunDialog element can be contained in a Maplet element; in an equation for the onapprove, oncancel, onchange, onclick, ondecline, or onstartup option for an element; or wrapped in an Action element as a parameter in an element that accepts an onchange or onclick option without an equation.\n • The following table describes the control and use of the RunDialog element options.\n An x in the I column indicates that the option can be initialized, that is, specified in the calling sequence (element definition).\n An x in the R column indicates that the option is required in the calling sequence.\n An x in the G column indicates that the option can be read, that is, retrieved by using the Get tool.\n An x in the S column indicates that the option can be written, that is, set by using the SetOption element or the Set tool.\n\n Option I R G S dialog x x\n\n • The opts argument can contain one or more of the following equations that set Maplet application options.\n dialog = reference to a dialog element (name or string)\n A reference to the dialog to run.\n\nExamples\n\nIn this Maplet application, the RunDialog element starts a dialog in the same Maplet application.\n\n > $\\mathrm{with}\\left(\\mathrm{Maplets}\\left[\\mathrm{Elements}\\right]\\right):$\n > $\\mathrm{maplet}≔\\mathrm{Maplet}\\left(\\mathrm{Window}\\left(\\left[\\left[\\mathrm{TextField}\\left['\\mathrm{TF1}'\\right]\\left(\\right)\\right],\\left[\\mathrm{Button}\\left(\"Differentiate with respect to x\",\\mathrm{Evaluate}\\left('\\mathrm{TF1}'='\\mathrm{diff}\\left(\\mathrm{TF1},x\\right)'\\right)\\right),\\mathrm{Button}\\left(\"Help\",\\mathrm{RunDialog}\\left('\\mathrm{MD1}'\\right)\\right),\\mathrm{Button}\\left(\"Exit\",\\mathrm{Shutdown}\\left(\\left['\\mathrm{TF1}'\\right]\\right)\\right)\\right]\\right]\\right),\\mathrm{MessageDialog}\\left['\\mathrm{MD1}'\\right]\\left(\"See ?diff for help with the differentiation command\"\\right)\\right):$\n > $\\mathrm{Maplets}\\left[\\mathrm{Display}\\right]\\left(\\mathrm{maplet}\\right)$\n\nIn this Maplet application, the Display function starts another Maplet application that contains a dialog.\n\n > $\\mathrm{with}\\left(\\mathrm{Maplets}\\left[\\mathrm{Elements}\\right]\\right):$\n > $\\mathrm{maplet1}≔\\mathrm{Maplet}\\left(\\mathrm{Window}\\left(\\left[\\left[\\mathrm{TextField}\\left['\\mathrm{TF1}'\\right]\\left(\\right)\\right],\\left[\\mathrm{Button}\\left(\"Differentiate with respect to x\",\\mathrm{Evaluate}\\left('\\mathrm{TF1}'='\\mathrm{diff}\\left(\\mathrm{TF1},x\\right)'\\right)\\right),\\mathrm{Button}\\left(\"Help\",\\mathrm{Evaluate}\\left(\\mathrm{function}='\\mathrm{Maplets}\\left[\\mathrm{Display}\\right]\\left(\\mathrm{maplet2}\\right)'\\right)\\right),\\mathrm{Button}\\left(\"Exit\",\\mathrm{Shutdown}\\left(\\left['\\mathrm{TF1}'\\right]\\right)\\right)\\right]\\right]\\right)\\right):$\n > $\\mathrm{maplet2}≔\\mathrm{Maplet}\\left(\\mathrm{MessageDialog}\\left(\"See ?diff for help with the differentiation command\"\\right)\\right):$\n > $\\mathrm{Maplets}\\left[\\mathrm{Display}\\right]\\left(\\mathrm{maplet1}\\right)$" ]
[ null, "https://bat.bing.com/action/0", null ]
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https://chem.libretexts.org/Courses/Sacramento_City_College/SCC%3A_CHEM_300_-_Beginning_Chemistry/SCC%3A_CHEM_300_-_Beginning_Chemistry_(Alviar-Agnew)/03%3A_Matter_and_Energy/3.11%3A_Temperature_Changes-_Heat_Capacity
[ "# 3.11: Temperature Changes- Heat Capacity\n\nIf a swimming pool and wading pool, both full of water at the same temperature, were subjected to the same input of heat energy, the wading pool would certainly rise in temperature more quickly than the swimming pool. The heat capacity of an object depends both on its mass and its chemical composition. Because of its much larger mass, the swimming pool of water has a larger heat capacity than the wading pool.\n\n## Heat Capacity and Specific Heat\n\nDifferent substances respond to heat in different ways. If a metal chair sits in the bright sun on a hot day, it may become quite hot to the touch. An equal mass of water in the same sun will not become nearly as hot. We would say that water has a high heat capacity (the amount of heat required to raise the temperature of an object by $$1^\\text{o} \\text{C}$$). Water is very resistant to changes in temperature, while metals in general are not. The specific heat of a substance is the amount of energy required to raise the temperature of 1 gram of the substance by $$1^\\text{o} \\text{C}$$. The symbol for specific heat is $$c_p$$, with the $$p$$ subscript referring to the fact that specific heats are measured at constant pressure. The units for specific heat can either be joules per gram per degree $$\\left( \\text{J/g}^\\text{o} \\text{C} \\right)$$ or calories per gram per degree $$\\left( \\text{cal/g}^\\text{o} \\text{C} \\right)$$. This text will use $$\\text{J/g}^\\text{o} \\text{C}$$ for specific heat.\n\nNotice that water has a very high specific heat compared to most other substances. Water is commonly used as a coolant for machinery because it is able to absorb large quantities of heat (see table above). Coastal climates are much more moderate than inland climates because of the presence of the ocean. Water in lakes or oceans absorbs heat from the air on hot days and releases it back into the air on cool days.", null, "Figure $$\\PageIndex{1}$$: This power plant in West Virginia, like many others, is located next to a large lake so that the water from the lake can be used as a coolant. Cool water from the lake is pumped into the plant, while warmer water is pumped out of the plant and back into the lake.\n\n## Summary\n\n• Heat capacity is the amount of heat required to raise the temperature of an object by $$1^\\text{o} \\text{C}$$).\n• The specific heat of a substance is the amount of energy required to raise the temperature of 1 gram of the substance by $$1^\\text{o} \\text{C}$$." ]
[ null, "https://chem.libretexts.org/@api/deki/files/78936/CK12_Screenshot_17-4-1.png", null ]
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http://www.scielo.org.mx/scielo.php?script=sci_arttext&pid=S0035-001X2019000300268&lng=en&nrm=iso&tlng=en
[ "## Indicators\n\n•", null, "Cited by SciELO\n•", null, "Access statistics\n\n•", null, "Similars in SciELO\n\n## Print version ISSN 0035-001X\n\n### Rev. mex. fis. vol.65 n.3 México May./Jun. 2019  Epub Apr 30, 2020\n\n#### https://doi.org/10.31349/revmexfis.65.268\n\nResearch\n\nNear-field analysis and field transformation applied to a parabolic profile at 5 GHz\n\naInstituto Politécnico Nacional, Escuela Superior de Ingeniería Mecánica y Eléctrica. UPALM Edif. Z-4, Tercer piso, C.P. 07738 Ciudad de México.\n\nbUniversidad de Quintana Roo Campus Chetumal. Boulevard Bahía, Esq. Ignacio Comonfort S/N Edif. L, Col. Del Bosque, 77019. e-mail: [email protected]\n\nAbstract\n\nWe propose a method of near-field analysis and field transformation that is applied to a 1.5 m diameter parabolic reflector at a frequency of 5 GHz. An antenna of such dimensions requires at least a 75 m region to obtain its radiation pattern, and this represents a problem, here arises the necessity to make field transformations like the one presented in this work. Near field is modeled by means of Finite Difference Time Domain Method (FDTD) and current distribution is obtained using the discrete Pocklington equation. Radiation pattern is calculated applying the array factor for parabolic profiles. Results are compared with those obtained by CST Microwave Studio with a very good agreement.\n\nKeywords: Arrays; FDTD; Fraunhofer; modeling; parabolic profile; Rayleigh\n\nPACS: 03.50.De; 02.70.Bf; 41.20.Jb; 07.05.Tp; 01.50.H-\n\nResumen\n\nProponemos un método para el análisis en campo cercano y transformación de campo, que es aplicado a un reflector parabólico de 1.5 m de diámetro, a una frecuencia de 5 GHz, una antena de tales dimensiones requiere una región de al menos 75 m para obtener su diagrama de radiación, lo cual representa un problema, aquí es donde surge la necesidad de realizar transformaciones de campo como la presentada en este artículo. El campo cercano es modelado mediante el método de Diferencias Finitas en el Dominio del Tiempo (MDFDT) y la distribución de corriente es obtenida usando la ecuación discreta de Pocklington. El diagrama de radiación es calculado aplicando el factor de arreglo para perfiles parabólicos. Nuestros resultados del análisis en campo cercano y transformación de campo, son comparados contra el modelado de la estructura en CST Microwave Studio, encontrando muy buenas coincidencias entre ambos métodos.\n\nDescriptores: Arreglos; MDFDT; Fraunhofer; modelado; perfil parabólico; Rayleigh\n\n1. Introduction\n\nWhen a radiated field is distributed in space, there are three regions that identify the process : Rayleigh region, Fresnel region and Fraunhofer region; Rayleigh region refers to the individual point field, which means that the radiation generated at each point of the structure has not been mixed with the other. In Fresnel region, the generated individual fields begin to interact with each other, but their phase differences are significant, so the radiation pattern does not have the shape that must adopt in far-field. In Fraunhofer region, fields are added almost in phase and as they move away from the source, the differences begin to be further reduced, so this region is referred as the far-field region. Figure 1 describes the three regions of the field distribution process in space. Table I shows conditions to determine the field region distances, where D is the antenna diameter, R1 and R2 are Rayleigh and Fresnel regions and λ is the wavelength .\n\nTABLE I Field regions distances.\n\n Antenna Rayleigh Fresnel Fraunhofer dimension region Region region D>>λ R1<λ/2π λ2π\n\nAccording with Fig. 1 and Table I, the minimum distance to obtain the far-field of a 1.5 m diameter structure, working at 5 GHz, is at least 75 m, which represents a problem if we try to measure far-field inside a controlled environment (as in an anechoic chamber, making nearly unmanageable to address the problem). Our method proposes to perform a field modeling at the Rayleigh region, define virtual point currents over the surface through discrete Pocklington equation, and then obtain far-field using array theory; essentially a near-field to far-field transformation is carried out. Comparing with commercial software, we reduce computational resources and the proposed method has been used to transform near-field into far-field from real measurements .\n\nThe method requires knowing the electric field near to reflector surface; this data can be obtained from measurements as [3-6] or by field modeling in the Rayleigh region as in , for this purpose, we implement a model based on FDTD method in order to obtain the magnitude and phase of field in a small region, near to the reflector surface. Descriptions of FDTD, discrete Pocklington equation and parabolic array factor (PAF) are presented in following Subsecs. 1.1 and 1.2.\n\n1.1 Finite difference time domain method (FDTD)\n\nThe FDTD method allows analyzing the effects of electromagnetic propagation, transforming the differential Maxwell equations into finite difference equations that can be handled by computers [8,9]. FDTD solution involves: establish the analysis region and divide it into cells where Maxwell’s equations are approximated by equations in finite differences, considering all materials within the region, using their electric and magnetic characteristics given by their permittivity (ε), permeability (μ) and conductivity (σ); finally the iterative computing solution gives the point field over the region.\n\nModeled parabolic profile was performed in two dimensions, then propagation is restricted to the XY plane. We choose TM mode considering Ez field as the only transversal component in the XY propagating plane, therefore Ex=Ey=Hz=0. FDTD requires that the analysis region be surrounded by a free reflection computational boundary; we choose Perfectly Matched Layer model (PML) for our application as in [9,10]. Equations (1-4) show the present field components.\n\nEzx(i,j)=abEzx(i,j)+Δtεb[Hy(i,j)-Hy(i-1,j)Δx] (1)\n\nEzy(i,j)=abEzy(i,j)-Δtεb[Hx(i,j)-Hx(i,j-1)Δy] (2)\n\nHx(i,j)=cdHx(i,j)-Δtμd[(Ezx+Ezy)(i.j+1)Δy-(Ezx+Ezy)(i,j)Δy] (3)\n\nHy(i,j)=cdHy(i,j)+Δtμd[(Ezx+Ezy)(i+1,j)Δx-(Ezx+Ezy)(i.j)Δx] (4)\n\nWhere:\n\na=1-σΔt2ε,b=1+σΔt2ε,\n\nC=1-σ*Δt2μ,d=1+σ*Δt2μ,\n\nAs a primary source we use a standard WR187 waveguide with dimensions of 0.1×0.02215 m working at 5 GHz; as the experiment was performed in 2D, field distribution is given by Eq. (5), with i=j=0 at the origin and E0 constant.\n\nEz(i,j)n=E0sin(2πfnΔt) (5)\n\n1.2 Discrete Pocklington equation and parabolic array factor\n\nThe method considers that the generalized Pocklington equation is valid not only on the antenna surface, but also outside very near to it, as near as the Rayleigh region defines it. Hence, it is possible to perform a field modeling by FDTD or by measurement as in to acquire field points near the structure.\n\nThe current can be determined by transforming the generalized Pocklington in a discrete form as Eq. (6) calculating virtual currents In along the structure.\n\nEmi=-1jωεn=1N[R2(k2R2-1-jkR)Smsn'+(3+3jkR-k2R2)(RSm)(Rsn')]e-jkR4πR5In. (6)\n\nWhere:\n\nRmn=|r-r'|=[x(s)-x'(s')]2+[y(s)-y'(s')]2+[z(s)-z'(s')]2\n\nk=Wavenumber\n\nSolving Eq. (6) by the structure geometry, it can be constructed a matrix as (7), where ZMN represent impedances, IN represents a virtual current to be calculated and EM represents the modeled electric field. The matrix is solved by multiplying the electric field vector, obtained by the FDTD modeling, by the inverse of the impedance matrix, so vector current is obtained.\n\n*20cZ11Z1NZM1ZMN*20cI1IN=*20cE1EM (7)\n\nFinally, far-field is calculated using the array theory [12,13] under the assumption that calculated point currents IN in (7), represent an array of specific elements that generate the far-field as is shown in Fig. 2. Field radiated by a parabolic array of point sources is given by (8):\n\nPAF(θ)=E0n=1N|IN|ej(kRncos(θ-ϕ)+αn) (8)\n\nIn (8) E0 represents the reference field of the array point sources with unitary magnitude; the sum of the individual fields from each point antenna considers the phase current αn, due the distance and the position of each current point Rn defined by ϕ while θ represents the angle between each point and the far-field position, Fig. 2 describes Eq. (8).\n\n2. Near-field analysis and far-field transformation\n\nTable II shows the modeling parameters, where physical dimensions are defined in meters and cells, in order to apply our computational technique based on Pocklington equation and FDTD.\n\nTABLE II Modeling parameters at 5 GHz.\n\n Units wavelength Cells λ m 0.06 Cell λ/20 0.003 Calculus region Plate diameter 500 × 750 500 1.5 × 2.25 1.5 Vertex - focus 200 0.6 Vertex - edge 77 0.23 height (feeder) 7 0.02215 depth (feeder) 33 0.1 PML 24 0.07\n\nCalculus region dimensions are 1.5×2.25 m corresponding to 500×750 cells, considering the size of each cell as λ/20, even several authors define λ/10 is an acceptable size . Our results are compared with those of commercially available software CST, using same parameters for both experiments. Figure 3 shows the calculus region and electric field distribution for both computational methods. To have a better view of the parabolic profile, Fig. 3 (a) is shown in a smaller area defined by (90:335,95:645), where reflector vertex is located at (100,375) and its focus at (300,375). Figure 3 displays both incident and reflected fields, which visual comparison with CST shows a very good concordance.", null, "FIGURE 3 Near-field modeling at 5 GHz, a) FDTD and b) CST.\n\nThe electric field vector EM in Rayleigh region required in (6), is defined using field samples from 500 points in the FDTD matrix at 0.003 m, one cell away from reflector surface along the parabolic profile. Figure 4 shows EM vector in magnitude, real and imaginary parts. Field phase is shown in Fig. 5, and is obtained using distances from reflector surface to its focus, for each field point using parabola equation.\n\nNext step in the methodology is to apply Eqs. (6) and (7) to obtain current along the reflector surface. Equation (6) requires the use of parabolic geometry for distace Rmn and vectors sm and sn, that is, distances between each point of the reflector axis and each current point .\n\nFinally, the field distribution is transformed in virtual point currents along the plate surface that will be applied in the array factor equation for the parabolic profile (8). Figure 6, shows the magnitude, real and imaginary normalized parts for current distribution IN. We can notice a similar behavior with electric field displayed on Fig. 4, including the 180 phase change of imaginary part due the reflection. Figure 7 shows the calculated phase for each current point.\n\nOnce field distribution is transformed into virtual currents along the parabolic profile, magnitude and phase values are used in the parabolic array factor Eq. (8) to get the far-field pattern shown in Fig. 8. As it is shown, radiation pattern has a main lobe width of 2.5 which is consistent with the theory of parabolic reflectors , also the nearest sidelobes are observed at -17 dB, other lobes have amplitudes between -20dB and -55 dB.\n\n3. Comparison\n\nIn previous work [7,16], calculated radiation pattern was compared versus a measured radiation pattern however, actual analysis allows us to observe and compare electric field behavior with those calculated in CST to show its viability.\n\nFigure 9 shows a comparison between near-field behavior for both simulations (FDTD and CST) in magnitude. Figures 10 and 11 show an excellent agreement of the near-field in real and imaginary parts; finally, near-field phase is shown in Fig. 12, as seen there is a clear similarity between FDTD and CST curves, specially real and imaginary parts with small differences at the beginning and at end of curves.", null, "FIGURE 9 Magnitude Electric field comparison between FDTD and CST.", null, "FIGURE 10 Electric field real part comparison between FDTD and CST.", null, "FIGURE 11 Electric field imaginary part comparison between FDTD and CST.", null, "FIGURE 12 Electric field phase comparison between FDTD and CST.\n\nRadiation patterns comparison is shown in Fig. 13, black curve represents the pattern of our method and blue curve is the CST pattern. Similarity between both curves is evident; it can be seen that the main lobe width is about 2.5 for both patterns and small differences of 5 dB at the edges (0° o 30° and 150° to 180°).\n\n4. Conclusions\n\nWe presented a method to transform near-field (Rayleigh region) into far-field, without performing simulation in large computational regions, then reducing computing resources. Even we define the field in the Rayleigh region using the FDTD Method, an alternative proposed as future work is to measure the field near the reflector surface and obtain far-field with Pocklington equation and Array Theory. In addition, our method can be applied to other antenna geometries as lineal radiators or even linear arrays or some other type of reflectors.\n\nComparison between our methodology (FDTD-Pocklington-PAF) and commercial software shows that results are very similar not only for radiation pattern, but also for near-field behavior.\n\nAcknowledgments\n\nThe authors want to thank to the Instituto Politécnico Nacional and Consejo Nacional de Ciencia y Tecnología of México for their support.\n\nReferences\n\n1. J. Sosa-Pedroza, L. Carrion-Rivera, F. Martínez-Zúñiga and S. Peña-Ruiz, Una Propuesta para Transformar Mediciones de Campo Cercano en Campo Lejano. Simposio de Metrología CENAM, México (2014). [ Links ]\n\n2. Y. Huang, K, Boyle, Antennas from Theory to Practice. (John Wiley. UK, 2008), p. 109-112. [ Links ]\n\n3. M. Sierra-Castañer, S. Burgos. Fresnel Zone to Far Field Algorithm for Rapid Array Antenna Measurements. IEEE EUCAP (2011). [ Links ]\n\n4. R. Cornelius, T. Salmerón-Ruiz, F. Saccardi, L. Foged, D. Heberling, M. Sierra-Castañer, IEEE Antennas and Propagation Magazine 56 (2014) . [ Links ]\n\n5. F. D’Agostino et al., IEEE Antennas and Propagation Magazine 54 (2012). [ Links ]\n\n6. S. Costanzo, G. Di Massa, Microwave and Optical Technology Letters 48 (2006). [ Links ]\n\n7. J. Sosa-Pedroza , M. Enciso-Aguilar, S. Peña-Ruiz, A. Rodríguez-Sánchez, and E. Garduño-Nolasco, Microwave and Optical Technology Letters 58 (2016). [ Links ]\n\n8. J. Sosa-Pedroza et al., Approach. Electron. Lett. 47 (2011) 1308-1309. [ Links ]\n\n9. A. Taflove and S. C. Hagness, Computational Electrodynamics The Finite-Diference Time-Domine Method. (Artech House USA 2005) p. 54-75. [ Links ]\n\n10. M. Benavides et al., Rev. Mex. Fís. E 57 (2011) 25-31. [ Links ]\n\n11. J. Sosa-Pedroza , V. Barrera-Figueroa and J. López-Bonilla, Pocklington Equation and the Method of Moments Proc. Pakistan Acad. (2005). [ Links ]\n\n12. C. A. Balanis, Antenna Theory Analysis and Design. 3rd Ed. (John Willey, 2005) p. 284-313. [ Links ]\n\n13. L. Josefsson, P. Persson, Conformal Array Antenna Theory and Design. IEEE (Press Series on Electromagnetic Wave Theory. USA, 2006), p. 15-22. [ Links ]\n\n14. Jacob W.M. Baars, The Paraboloidal Reflector Antenna in Radio Astronomy and Communication. (Springer Science USA, 2007), p. 14-16. [ Links ]\n\n15. J. D. Kraus, R. J. Marhefka, Antennas for all Applications. 3rd Ed. (McGraw Hill. USA, 2002), p. 564-571. [ Links ]\n\n16. J. Sosa, S. Peña, F. Martínez, Parabolic Reflector Near-field to Far-field Transformation Using FDTDM and Pocklington Equation. PIERS, (2017). [ Links ]\n\nReceived: August 11, 2018; Accepted: January 14, 2018", null, "This is an open-access article distributed under the terms of the Creative Commons Attribution License" ]
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http://pressurevesseltech.asmedigitalcollection.asme.org/article.aspx?articleid=2666607&resultClick=1
[ "0\nResearch Papers: SPECIAL SECTION PAPERS\n\n# Analytical and Semi-Analytical Methods for the Evaluation of Dynamic Thermo-Elastic Behavior of Structures Resting on a Pasternak Foundation\n\n[+] Author and Article Information\nXu Liang, Zeng Cao, Xing Zha, Jianxing Leng\n\nOcean College,\nZhejiang University,\nHangzhou 310058, Zhejiang, China\n\nHongyue Sun\n\nOcean College,\nZhejiang University,\nHangzhou 310058, Zhejiang, China\ne-mail: [email protected]\n\n1Corresponding author.\n\nContributed by the Pressure Vessel and Piping Division of ASME for publication in the JOURNAL OF PRESSURE VESSEL TECHNOLOGY. Manuscript received August 9, 2017; final manuscript received December 9, 2017; published online December 14, 2018. Assoc. Editor: Fabrizio Paolacci.\n\nJ. Pressure Vessel Technol 141(1), 010908 (Dec 14, 2018) (10 pages) Paper No: PVT-17-1146; doi: 10.1115/1.4038724 History: Received August 09, 2017; Revised December 09, 2017\n\n## Abstract\n\nAn analytical method and a semi-analytical method are proposed to analyze the dynamic thermo-elastic behavior of structures resting on a Pasternak foundation. The analytical method employs a finite Fourier integral transform and its inversion, as well as a Laplace transform and its numerical inversion. The semi-analytical method employs the state space method, the differential quadrature method (DQM), and the numerical inversion of the Laplace transform. To demonstrate the two methods, a simply supported Euler–Bernoulli beam of variable length is considered. The governing equations of the beam are derived using Hamilton's principle. A comparison between the results of analytical method and the results of semi-analytical method is carried out, and it is shown that the results of the two methods generally agree with each other, sometimes almost perfectly. A comparison of natural frequencies between the semi-analytical method and the experimental data from relevant literature shows good agreements between the two kinds of results, and the semi-analytical method is validated. Different numbers of sampling points along the axial direction are used to carry out convergence study. It is found that the semi-analytical method converges rapidly. The effects of different beam lengths and heights, thermal stress, and the spring and shear coefficients of the Pasternak medium are also investigated. The results obtained in this paper can serve as benchmark in further research.\n\n<>\n\n## References\n\nBoley, B. A. , 1956, “ Thermally Induced Vibrations of Beams,” J. Aeronaut. Sci., 23(2), pp. 179–181.\nAvsec, J. , and Oblak, M. , 2007, “ Thermal Vibrational Analysis for Simply Supported Beam and Clamped Beam,” J. Sound Vib., 308(3), pp. 514–525.\nMurmu, T. , and Pradhan, S. C. , 2009, “ Thermo-Mechanical Vibration of a Single-Walled Carbon Nanotube Embedded in an Elastic Medium Based on Nonlocal Elasticity Theory,” Comput. Mater. Sci., 46(4), pp. 854–859.\nPradhan, S. C. , and Murmu, T. , 2009, “ Thermo-Mechanical Vibration of FGM Sandwich Beam Under Variable Elastic Foundations Using Differential Quadrature Method,” J. Sound Vib., 321(1–2), pp. 342–362.\nJandaghian, A. A. , and Rahmani, O. , 2016, “ Free Vibration Analysis of Magneto-Electro-Thermo-Elastic Nanobeams Resting on a Pasternak Foundation,” Smart Mater. Struct., 25(3), p. 035023.\nLiang, F. , Li, Y. , and Huang, M. , 2013, “ Simplified Method for Laterally Loaded Piles Based on Pasternak Double-Parameter Spring Model for Foundations,” Chin. J. Geotech. Eng., 35(Suppl. 1), pp. 300–304.\nSharma, J. N. , and Kaur, R. , 2015, “ Modeling and Analysis of Forced Vibrations in Transversely Isotropic Thermoelastic Thin Beams,” Meccanica, 50(1), pp. 189–205.\nSharma, J. N. , and Kaur, R. , 2014, “ Analysis of Forced Vibrations in Micro-Scale Anisotropic Thermo-Elastic Beams Due to Concentrated Loads,” J. Therm. Stresses, 37(1), pp. 93–116.\nSong, X. , and Li, S.-R. , 2007, “ Thermal Buckling and Post-Buckling of Pinned–Fixed Euler–Bernoulli Beams on an Elastic Foundation,” Mech. Res. Commun., 34(2), pp. 164–171.\nNejad, M. Z. , Hadi, A. , and Rastgoo, A. , 2016, “ Buckling Analysis of Arbitrary Two-Directional Functionally Graded Euler–Bernoulli Nano-Beams Based on Nonlocal Elasticity Theory,” Int. J. Eng. Sci., 103, pp. 1–10.\nLi, S. , and Song, X. , 2006, “ Large Thermal Deflections of Timoshenko Beams Under Transversely Non-Uniform Temperature Rise,” Mech. Res. Commun., 33(1), pp. 84–92.\nYun, Y. , Chen, J. , Zhao, K. , Yan, B. , and Cao, H. , 2014, “ Analysis of Thermo-Elastic Coupling Effects of the Beam Structure With Interval Parameters,” J. Xidian Univ., 41(4), pp. 64–70.\nShen, Z. , Tian, Q. , Liu, X. , and Hu, G. , 2013, “ Thermally Induced Vibrations of Flexible Beams Using Absolute Nodal Coordinate Formulation,” Aerosp. Sci. Technol., 29(1), pp. 386–393.\nBert, C. W. , and Malik, M. , 1996, “ Differential Quadrature Method in Computational Mechanics: A Review,” ASME Appl. Mech. Rev., 49(1), pp. 1–28.\nLiang, X. , Wu, Z. , Wang, L. , Liu, G. , Wang, Z. , and Zhang, W. , 2014, “ Semianalytical Three-Dimensional Solutions for the Transient Response of Functionally Graded Material Rectangular Plates,” J. Eng. Mech., 141(9), p. 04015027.\nAkbarzadeh, A. H. , Hosseini, K. , Eslami, M. R. , and Sadighi, M. , 2011, “ Mechanical Behaviour of Functionally Graded Plates Under Static and Dynamic Loading,” Proc. Inst. Mech. Eng., Part C, 225(2), pp. 326–333.\nKiani, Y. , Shakeri, M. , and Eslami, M. R. , 2012, “ Thermoelastic Free Vibration and Dynamic Behaviour of an FGM Doubly Curved Panel Via the Analytical Hybrid Laplace–Fourier Transformation,” Acta Mech., 223(6), pp. 1199–1218.\nAkbarzadeh, A. H. , Abbasi, M. , and Eslami, M. R. , 2012, “ Coupled Thermoelasticity of Functionally Graded Plates Based on the Third-Order Shear Deformation Theory,” Thin-Walled Struct., 53, pp. 141–155.\nWang, Z. , Liang, X. , and Liu, G. , 2013, “ An Analytical Method for Evaluating the Dynamic Response of Plates Subjected to Underwater Shock Employing Mindlin Plate Theory and Laplace Transforms,” Math. Probl. Eng., 2013, pp. 927–940.\nDurbin, F. , 1974, “ Numerical Inversion of Laplace Transforms: An Efficient Improvement to Dubner and Abate's Method,” Comput. J., 17(4), pp. 371–376.\nLiang, X. , Wang, Z. , Wang, L. , and Liu, G. , 2014, “ Semi-Analytical Solution for Three-Dimensional Transient Response of Functionally Graded Annular Plate on a Two Parameter Viscoelastic Foundation,” J. Sound Vib., 333(12), pp. 2649–2663.\nLiang, X. , Kou, H. L. , Wang, L. , Palmer, A. C. , Wang, Z. , and Liu, G. , 2015, “ Three-Dimensional Transient Analysis of Functionally Graded Material Annular Sector Plate Under Various Boundary Conditions,” Compos. Struct., 132, pp. 584–596.\nZavodney, L. D. , and Nayfeh, A. H. , 1988, “ The Response of a Single-Degree-of-Freedom System With Quadratic and Cubic Non-Linearities to a Fundamental Parametric Resonance,” J. Sound Vib., 120(1), pp. 63–93.\nShu, C. , and Du, H. , 1997, “ Implementation of Clamped and Simply Supported Boundary Conditions in the GDQ Free Vibration Analysis of Beams and Plates,” Int. J. Solids Struct., 34(7), pp. 819–835.\nWolfram Research, 2016, “Wolfram Mathematica, Version 11.0.0,” Wolfram Research, Inc., Champaign, IL.\nMcLellan, R. B. , and Ishikawa, T. , 1987, “ The Elastic Properties of Aluminum at High Temperatures,” J. Phys. Chem. Solids, 48(7), pp. 603–606.\nNix, F. C. , and Macnair, D. , 1941, “ The Thermal Expansion of Pure Metals: Copper, Gold, Aluminum, Nickel, and Iron,” Phys. Rev., 60(8), pp. 597–605.\nMarques, R. F. A. , and Inman, D. J. , 2002, “An Analytical Model for a Clamped Isotropic Beam Under Thermal Effects,” ASME Paper No. IMECE2002-33977.\n\n## Figures\n\nFig. 1\n\nThe geometry of a thermo-elastic beam on a Pasternak foundation", null, "Fig. 2\n\nThe time histories of peak beam deflection at the selected position, x = l/2, in the thermal environment by the analytical and the semi-analytical method: (a) case 1, (b) case 2, (c) case 3, and (d) case 4", null, "Fig. 3\n\nThe time histories of peak beam deflection at a selected point, x = l/2, of the beam, for four different numbers (N) of sampling points along the x-axis: (a) case 1, (b) case 2, (c) case 3, and (d) case 4", null, "Fig. 4\n\nThe time histories of peak beam deflection at two selected positions, x = l/2 and x = l/5, for different boundary conditions: (a) at the point x = l/2 and (b) at the point x = l/5", null, "Fig. 5\n\nThe time histories of peak beam deflection at two selected positions, x = l/2 and x = l/5, for different beam lengths and fixed height: (a) at the point x = l/2 and (b) at the point x = l/5", null, "Fig. 6\n\nThe time histories of peak beam deflection at two selected points, x = l/2 and x = l/5, for a given beam length and different heights: (a) at the point x = l/2 and (b) at the point x = l/5", null, "Fig. 7\n\nThe time histories of peak beam deflection at two different selected points, x = l/2 and x = l/5, for different surrounding temperatures: (a) at the point x = l/2 and (b) at the point x = l/5", null, "Fig. 8\n\nThe time histories of peak beam deflection at two different selected points, x = l/2 and x = l/5, with different spring coefficients: (a) at the point x = l/2 and (b) at the point x = l/5", null, "Fig. 9\n\nThe time histories of peak beam deflection at two different selected points x = l/2 and x = l/5 with different shear coefficients: (a) at the point x = l/2 and (b) at the point x = l/5", null, "## Errata\n\nSome tools below are only available to our subscribers or users with an online account.\n\n### Related Content\n\nCustomize your page view by dragging and repositioning the boxes below.\n\nRelated Journal Articles\nRelated Proceedings Articles\nRelated eBook Content\nTopic Collections" ]
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https://www.tutorialspoint.com/algorithm-to-sum-ranges-that-lie-within-another-separate-range-in-javascript
[ "# Algorithm to sum ranges that lie within another separate range in JavaScript\n\nJavascriptWeb DevelopmentFront End TechnologyObject Oriented Programming\n\nWe have two sets of ranges; one is a single range of any length (R1) and the other is a set of ranges (R2) some or parts of which may or may not lie within the single range (R1).\n\nWe need to calculate the sum of the ranges in (R2) - whole or partial - that lie within the single range (R1).\n\nconst R1 = [20,40];\nconst R2 = [[14,22],[24,27],[31,35],[38,56]];\n\nResult\n\n= 2+3+4+2 = 11\nR1 = [120,356];\nR2 = [[234,567]];\n\nResult\n\n122\n\n## Example\n\nLet us write the code −\n\nconst R1 = [20,40];\nconst R2 = [[14,22],[24,27],[31,35],[38,56]];\nconst R3 = [120,356];\nconst R4 = [[234,567]];\nfunction sumRanges(range, values) {\nconst [start, end] = range;\nconst res = values.reduce((acc, val) => {\nconst [left, right] = val;\nconst ex1 = Math.min(right, end);\nconst ex2 = Math.max(left, start);\nconst diff = ex1 - ex2;\nreturn acc + Math.max(0, diff);\n}, 0);\nreturn res;\n};\nconsole.log(sumRanges(R1, R2));\nconsole.log(sumRanges(R3, R4));\n\n## Output\n\nAnd the output in the console will be −\n\n11\n122" ]
[ null ]
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https://www.physicsforums.com/threads/double-integrals.312276/
[ "# Double integrals\n\n## Homework Statement\n\nintegrate y DA over the regions s, where s is in the first quadrant bounded by x = y^2 and x = 8 - y^2\n\n## The Attempt at a Solution\n\nIf the y wasnt there i would evaluate two integrals and subtract them to get the area.\nBut since the y is there, can i still evaluate them seperatly with the y in place and subtract.\n\nRelated Calculus and Beyond Homework Help News on Phys.org\ntiny-tim\nHomework Helper\nHi joemama69!", null, "integrate y DA over the regions s, where s is in the first quadrant bounded by x = y^2 and x = 8 - y^2\nand which axis??\n\nIf the y wasnt there i would evaluate two integrals and subtract them to get the area.\nBut since the y is there, can i still evaluate them seperatly with the y in place and subtract.\n(not sure what you're subtracting, but:) yes … what's worrying you about that?", null, "ya i was thinkin to hard\n\nheres what i did, i broke it into two intigrals and added them\n\nfirst intigral\n\ni differentiated y dy from 0 to x^(1/2)\nthen i did in trems of x from 0 to 4 and got an area for 4\n\n2nd integral\n\ni differentiated y dy from 0 to (8-x)^(1/2)\nthen i did in trems of x from 4 to 8 and got an area for 4\n\n4 + 4 = 8 is that right\n\ntiny-tim\nHomework Helper\nya i was thinkin to hard\n\nheres what i did, i broke it into two intigrals and added them\n\n4 + 4 = 8 is that right\nYup, that's fine!", null, "(though you could have saved a little work by pointing out that the region was obviously symmetric about x = 4, so all you had to do was evaluate one integral, and then double it.", null, ")\n\nya i noticed that, but i figured for my teachers sake i should do it thoroughly" ]
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https://electronics.stackexchange.com/questions/363601/issue-with-simulink-does-not-support-code-generation
[ "# Issue with “Simulink does not support code generation”\n\nI have something similar to this simplified following block diagram on Simulink. It looks rather messy.", null, "I want to replace a 3-point summing block with a function block, while keeping the same output.", null, "First I started by placing the code inside the function block:\n\nfunction y = fcn(u)\nsys1 = tf(0.5,[1 0 0 4]);\nsys2 = tf([3 0.5],[1 0 15]);\nsys3 = tf(1,[1 1]);\ny = sys1 + sys2 + sys3;\n\n\nHowever I was greeted with an error saying that Simulink does not support code generation.\n\n\"The 'tf' class does not support code generation.\"\n\nI am trying to implement an extrinsic function or'wrapper function' with some difficulty. I created a new script called myWrapper.m, containing the same code:\n\nfunction y = myWrapper(u)\nsys1 = tf(0.5,[1 0 0 0 4]);\nsys2 = tf([3 5],[1 0 15]);\nsys3 = tf(1,[1 1]);\ny = sys1 + sys2 + sys3;\n\n\nand the MATLAB Function edited to:\n\nfunction y1 = fcn(u1)\n\ny1 = myWrapper(u1);\n\n\nThe error persists.\n\nI somehow want to access myWrapper.m file from the MATLAB Function block. Any pointers on how this should be done? Following the previous link given and the official docs I am ending up with something like this in my MATLAB Function block:\n\nfunction y = fcn(u)coder.extrinsic('myWrapper')\n\ny = myWrapper(u);\n\n\nThe last code above is syntactically incorrect and I am at a loss on how it should be done. MATLAB automaticaly corrects the above code to:\n\nfunction y = fcn(u,coder,extrinsic, myWrapper )\n\ny = myWrapper(u);\n\n\nwhich is not what I want.\n\nAny tips and/or suggestions on how this could be done would be appreciated.\n\nA similar question was asked on the MathWorks forum here, two years ago, with no response.\n\n• This question would be better asked on the Mathworks forums. Short answer, not all MATLAB command are supported inside MATLAB function blocks. – scotty3785 Mar 20 '18 at 11:40\n• @scotty3785, thanks for your contribution. I asked the question over there and no answer. A similar question was asked on the Mathworks forums some years ago, and once again, no response. – Rrz0 Mar 20 '18 at 11:41\n\nI was going about tackling this problem completely wrong. In order to replace the summing block, one must NOT remove the Transfer Function blocks which feed into the summing block and add them inside the script.\n\nA MATLAB Function does not support code generation (and rightly so) such that a transfer function may be implemented inside it. That is why the blocks simply feed into the MATLAB Function as follows.", null, "The script would very simply be:\n\nfunction y1 = fcn(u1, u2, u3)\n\nx = (u1 + u2 +u3);\ny1 = x;\n\nend" ]
[ null, "https://i.stack.imgur.com/mOYZP.png", null, "https://i.stack.imgur.com/5sFWo.png", null, "https://i.stack.imgur.com/ZB4bg.png", null ]
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https://statsidea.com/normalize-information-in-sas/
[ "# Normalize Information in SAS\n\nTo “normalize” a collection of information values approach to scale the values such that the ruthless of the entire values is 0 and the usual bypass is 1.\n\nThis educational explains how you can normalize knowledge in SAS.\n\n### Instance: Normalize Information in SAS\n\nThink we now have please see dataset:", null, "Carry out please see steps to normalize this eager of information values in SAS.\n\nStep 1: Form the Dataset\n\nFirst, let’s virtue please see code to build the dataset in SAS:\n\n```/*build dataset*/\nknowledge original_data;\nenter values;\ndatalines;\n12\n14\n15\n15\n16\n17\n18\n20\n24\n25\n26\n29\n32\n34\n37\n;\nrun;\n\n/*view ruthless and same old bypass of dataset*/\nproc approach knowledge=original_data Ruthless StdDev ndec=3;\nvar values;\nrun;```", null, "From the output we will see that the ruthless of the dataset is 22.267 and the usual bypass is 7.968.\n\nStep 2: Normalize the Dataset\n\nLater, we’ll virtue proc stdize to normalize the dataset:\n\n```/*normalize the dataset*/\nproc stdize knowledge=original_data out=normalized_data;\nvar values;\nrun;\n\n/*print normalized dataset*/\nproc print knowledge=normalized_data;\n\n/*view ruthless and same old bypass of normalized dataset*/\nproc approach knowledge=normalized_data Ruthless StdDev ndec=2;\nvar values;\nrun;```", null, "From the output we will see that the ruthless of the normalized dataset is 0 and the usual bypass is 1.\n\nStep 3: Interpret the Normalized Information\n\nSAS old please see formulation to normalize the knowledge values:\n\nNormalized worth = (x – x) / s\n\nthe place:\n\n• x = knowledge worth\n• x = ruthless of dataset\n• s = same old bypass of dataset\n\nEach and every normalized worth tells us what number of same old deviations the actual knowledge worth used to be from the ruthless.\n\nAs an example, imagine the knowledge level “12” in our actual dataset. The actual pattern ruthless used to be 22.267 and the actual pattern same old bypass used to be 7.968.\n\nThe normalized worth for “12” became out to be -1.288, which used to be calculated as:\n\nNormalized worth = (x – x) / s = (12 – 22.267) / 7.968 = -1.288\n\nThis tells us that the price “12” is 1.288 same old deviations under the ruthless within the actual dataset.\n\nEach and every of the normalized values within the dataset can backup us know the way alike or some distance a specific knowledge worth is from the ruthless.\n\nA tiny normalized worth signifies {that a} worth is alike to the ruthless occasion a massive normalized worth signifies {that a} worth is some distance from the ruthless.\n\n### Extra Assets\n\nRefer to tutorials provide an explanation for how you can carry out alternative usual duties in SAS:\n\nSignificance Proc Abstract in SAS\nCalculate Correlation in SAS\nForm Frequency Tables in SAS" ]
[ null, "https://www.statology.org/wp-content/uploads/2021/12/norm1.jpg", null, "https://www.statology.org/wp-content/uploads/2021/12/norm3.jpg", null, "https://www.statology.org/wp-content/uploads/2021/12/norm2.jpg", null ]
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https://www.bartleby.com/questions-and-answers/water-flows-uniformly-in-a-trapezoidal-channel-with-a-14ftwide-bed-and-2.21-m-2.2-side-slopes.-if-th/90764a28-96b6-470f-bfee-7e16700505ef
[ "# Water flows uniformly in a trapezoidal channel with a 14-ft-wide bed and 2.2:1 (m = 2.2) side slopes. If the bed slope is 1.1 ft/mile and n = 0.015, find the flow rate when the depth is 9 ft.\n\nQuestion\n\nWater flows uniformly in a trapezoidal channel with a 14-ft-wide bed and 2.2:1 (m = 2.2) side slopes. If the bed slope is 1.1 ft/mile and n = 0.015, find the flow rate when the depth is 9 ft.\n\n### Want to see this answer and more?\n\nExperts are waiting 24/7 to provide step-by-step solutions in as fast as 30 minutes!*", null, "" ]
[ null, "https://www.bartleby.com/static/bartleby-logo-tm-tag-inverted.svg", null ]
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http://arguinea.info/block-wall-calculation/block-wall-calculation-average-concrete-masonry-units-and-mortar-per-sq-ft-of-wall-concrete-block-wall-calculator/
[ "# Block Wall Calculation Average Concrete Masonry Units And Mortar Per Sq Ft Of Wall Concrete Block Wall Calculator", null, "block wall calculation average concrete masonry units and mortar per sq ft of wall concrete block wall calculator.\n\nblock wall calculator uk rates analysis for calculating material and labour building works concrete load calculation,retaining wall calculator how to calculate block for cost formula estimator,block wall calculator uk figure square footage for flooring hardwood floor how to calculate retaining design calculations wind load calculation,half brick walling calculations block wall calculator cost retaining design concrete example uk,concrete block wall estimating tools for retaining walls courtyard and fences calculator design example uk,arch heat loss assignment winter concrete block retaining wall design example calculator uk calculations,block wall calculations detail concrete cost calculator nz retaining calculation,concrete block retaining wall design example how to calculate no of blocks in a building calc calculator metric,concrete block retaining wall design example uk calc tools mutual materials formula,block wall calculator south africa metric final building construction 1 australia." ]
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https://www.futuristicmath.com/math-quiz-games.html
[ "# Math Quiz Games For Children Online\n\nMath Quiz Games For Children Online.This page is a collection of math tests based on a jungle quiz theme. Children are asked to solve a problem and choose the correct answer. The quizzes are interactive in the sense that players can track their scores on the go. Its a free quiz hence players can practice until the skill is mastered. To take any quiz below, click on the play now button or on the image. You will be taken to the page where the game loads. Play the game until you master the skill being taight. In the end print out worksheets and get more practice. These activities are free and always online. Have fun playing.", null, "", null, "Addition of numbers twice e.g 4 + 4 = ?", null, "Addition of three digit numbers game", null, "Learn the principle of addition of zero.\n\n### Division Game", null, "Division of numbers board game\n\n### Even & Odd Numbers", null, "Even and odd numbers math board game\n\n### Multiply by single digits", null, "Learn basic multiplication of numbers game\n\n### Multiplication by ten game", null, "Board game on multiplication online\n\n### Multiply by five digits", null, "Board game on multiplication by five digits - Game based learning for kids\n\n### Number spelling game > 20", null, "A math board game on how to spell numbers. Review math concepts online\n\n### Numbers spelling up to 20", null, "Teach kids spelling numbers in a fun way. Free learning tools for parents.\n\n### Order of operations game", null, "Practice solving order of operations with fun - Free learning resources for teachers.\n\n### Multiply by five digits", null, "Board game on multiplication by five digits - Pirate gamer video game", null, "A math board game on pre-algebra. -Pirate game today\n\n### Pre-algebra subtraction", null, "Pre-algebra subtraction game online - Pirate game score\n\n### Prime & Composite numbers", null, "A game on identifying both number types - Game pirate of the Caribbean..\n\n### Roman numerals", null, "Board game on Roman to Arabic numerals - Pirate game for children\n\n### Subtraction - 3 digits", null, "A math board game on subtraction for students - Pirate Game on PC\n\n### Two digit subtraction", null, "Teach kids how to subtract two digit numbers. Pirate Game for review\n\n### Subtraction up to 20", null, "Practice subtraction up to twenty game - Interactive online learning" ]
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https://www.crackverbal.com/solutions/value-antique-watch/
[ "# Solutions\n\nGet detailed explanations to advanced GMAT questions.\n\n### Question\n\nAt the end of each year, the value of a certain antique watch is “c” percent more than its value one year earlier, where “c” has the same value each year. If the value of the watch was “k” dollars on January 1, 1992, and “m” dollars on January 1, 1994, then in terms of “m” and “k”, what was the value of the watch, in dollars, on January 1, 1995?\n\nOption A:\n\nm+1/2(m-k) B. m+1/2((m-k)/k)m C. (m*sqrt(m))/sqrt(k) D. m^2/2k; E. km^2\n\nOption B:\n\nm+1/2((m-k)/k)m\n\nOption C:\n\n(m*sqrt(m))/sqrt(k)\n\nOption D:\n\nm^2/2k\n\nOption E:\n\nkm^2\n\nHard\n\n### Solution\n\nOption C is the correct answer.\n\n### Option Analysis\n\nPrice in 1992 – k;\n\nPrice in 1993 – k∗(1+(c/100));\n\nPrice in 1994 – k∗(1+(c/100))^2=m –> (1+(c/100))=√(m/k);\n\nPrice in 1995 – m*(1+(c/100)) = m *√(m/k)" ]
[ null ]
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https://codereview.stackexchange.com/questions/164129/haskell-adjacency-matrix
[ "This program is used to solve Challenge #140 [Intermediate] on /r/dailyprogrammer. The challenge is, given a input set of directed edges, print the corresponding adjacency matrix.\n\nHere is a sample input:\n\n5 5\n0 3 -> 1\n1 -> 2\n2 -> 4\n3 -> 4\n0 -> 0 3\n\n\nand the corresponding solution:\n\n11010\n00100\n00001\n01001\n00000\n\n\nNote that there can be more than one edge per line - the set of edges generated is then the Cartesian product of those nodes. For instance, 0 3 -> 1 would generate the edges [(0, 1), (3, 1)].\n\nHere is my solution:\n\nimport Data.Function\nimport Data.Ix\nimport Data.List\nimport qualified Data.Set as Set\n\ntype Edge = (Integer, Integer)\n\ngenGrid :: Integer -> Set.Set Edge -> [[Bool]]\ngenGrid n edges = (map . map) (Set.member edges) coordGrid\nwhere\ncoordGrid = groupBy ((==) on fst) coords\ncoords = range ((0, 0), (n - 1, n - 1))\n\nparseEdge :: String -> [Edge]\nparseEdge s = [(from, to) | from <- outgoing, to <- incoming]\nwhere\noutgoing = map read (takeWhile (/= \"->\") split)\nincoming = map read (tail $dropWhile (/= \"->\") split) split = words s doChallenge :: [String] -> [String] doChallenge s = (map . map) pprint (genGrid n edges) where n = read (head$ words $head s) edges = Set.fromList$ concatMap parseEdge (drop 1 s)\npprint True = '1'\npprint _ = '0'\n\ninput :: [String]\ninput = [\"5 5\", \"0 3 -> 1\", \"1 -> 2\", \"2 -> 4\", \"3 -> 4\", \"0 -> 0 3\"]\n\noutput :: [String]\noutput = [\"11010\", \"00100\", \"00001\", \"01001\", \"00000\"]\n\nmain :: IO ()\nmain = print \\$ output == doChallenge input\n\n\nMy main points of focus" ]
[ null ]
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https://www.knowdebt.org/vocabulary/potential/
[ "In linguistics, the potential mood The mathematical study of potentials is known as potential theory; it is the study of harmonic functions on manifolds. This mathematical formulation arises from the fact that, in physics, the scalar potential is irrotational, and thus has a vanishing Laplacian — the very definition of a harmonic function. In physics, a potential may refer to the scalar potential or to the vector potential." ]
[ null ]
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https://demo.formulasearchengine.com/wiki/Intersection_theory
[ "# Intersection theory\n\nNot to be confused with Intersectionality theory.\n\nIn mathematics, intersection theory is a branch of algebraic geometry, where subvarieties are intersected on an algebraic variety, and of algebraic topology, where intersections are computed within the cohomology ring. The theory for varieties is older, with roots in Bézout's theorem on curves and elimination theory. On the other hand the topological theory more quickly reached a definitive form.\n\n## Topological intersection form\n\n{{#invoke:see also|seealso}} For a connected oriented manifold Template:Mvar of dimension 2n the intersection form is defined on the Template:Mvar-th cohomology group (what is usually called the 'middle dimension') by the evaluation of the cup product on the fundamental class [M] in H2n(M, ∂M). Stated precisely, there is a bilinear form\n\n$\\lambda _{M}\\colon H^{n}(M,\\partial M)\\times H^{n}(M,\\partial M)\\to \\mathbf {Z}$", null, "given by\n\n$\\lambda _{M}(a,b)=\\langle a\\smile b,[M]\\rangle \\in \\mathbf {Z}$", null, "with\n\n$\\lambda _{M}(a,b)=(-1)^{n}\\lambda _{M}(b,a)\\in \\mathbf {Z} .$", null, "This is a symmetric form for Template:Mvar even (so 2n = 4k doubly even), in which case the signature of Template:Mvar is defined to be the signature of the form, and an alternating form for Template:Mvar odd (so 2n= 4k + 2 singly even). These can be referred to uniformly as ε-symmetric forms, where ε = (−1)n = ±1 respectively for symmetric and skew-symmetric forms. It is possible in some circumstances to refine this form to an [[ε-quadratic form|Template:Mvar-quadratic form]], though this requires additional data such as a framing of the tangent bundle. It is possible to drop the orientability condition and work with Z/2Z coefficients instead.\n\nThese forms are important topological invariants. For example, a theorem of Michael Freedman states that simply connected compact 4-manifolds are (almost) determined by their intersection forms up to homeomorphism – see intersection form (4-manifold).\n\nBy Poincaré duality, it turns out that there is a way to think of this geometrically. If possible, choose representative Template:Mvar-dimensional submanifolds Template:Mvar, Template:Mvar for the Poincaré duals of Template:Mvar and Template:Mvar. Then λM (a, b) is the oriented intersection number of Template:Mvar and Template:Mvar, which is well-defined because of the dimensions of Template:Mvar and Template:Mvar.Template:Clarify This explains the terminology intersection form.\n\n## Intersection theory in algebraic geometry\n\nWilliam Fulton in Intersection Theory (1984) writes\n\n... if Template:Mvar and Template:Mvar are subvarieties of a non-singular variety Template:Mvar, the intersection product A · B should be an equivalence class of algebraic cycles closely related to the geometry of how AB, Template:Mvar and Template:Mvar are situated in Template:Mvar. Two extreme cases have been most familiar. If the intersection is proper, i.e. dim(AB) = dim A + dim B − dim X, then A · B is a linear combination of the irreducible components of AB, with coefficients the intersection multiplicities. At the other extreme, if A = B is a non-singular subvariety, the self-intersection formula says that A · B is represented by the top Chern class of the normal bundle of Template:Mvar in Template:Mvar.\n\nTo give a definition, in the general case, of the intersection multiplicity was the major concern of André Weil's 1946 book Foundations of Algebraic Geometry. Work in the 1920s of B. L. van der Waerden had already addressed the question; in the Italian school of algebraic geometry the ideas were well known, but foundational questions were not addressed in the same spirit.\n\n### Moving cycles\n\nA well-working machinery of intersecting algebraic cycle Template:Mvar and Template:Mvar requires more than taking just the set-theoretic intersection of the cycles in question. Certainly, the intersection VW or, more commonly called intersection product, denoted V · W, should consist of the set-theoretic intersection of the two subvarieties. However it occurs that cycles are in bad position, e.g. two parallel lines in the plane, or a plane containing a line (intersecting in 3-space). In both cases the intersection should be a point, because, again, if one cycle is moved, this would be the intersection. The intersection of two cycles Template:Mvar and Template:Mvar is called proper if the codimension of the (set-theoretic) intersection VW is the sum of the codimensions of Template:Mvar and Template:Mvar, respectively, i.e. the \"expected\" value.\n\nTherefore the concept of moving cycles using appropriate equivalence relations on algebraic cycles is used. The equivalence must be broad enough that given any two cycles Template:Mvar and Template:Mvar, there are equivalent cycles V′ and W′ such that the intersection V′W′ is proper. Of course, on the other hand, for a second equivalent V′′ and W′′, V′W′ needs to be equivalent to V′′W′′.\n\nFor the purposes of intersection theory, rational equivalence is the most important one. Briefly, two Template:Mvar-dimensional cycles on a variety Template:Mvar are rationally equivalent if there is a rational function f on a (k + 1)-dimensional subvariety Template:Mvar, i.e. an element of the function field k(Y) or equivalently a function f  : YP1, such that VW =  f−1(0) −  f−1(∞), where f−1(⋅) is counted with multiplicities. Rational equivalence accomplishes the needs sketched above.\n\n### Intersection multiplicities\n\nThe guiding principle in the definition of intersection multiplicities of cycles is continuity in a certain sense. Consider the following elementary example: the intersection of a parabola y = x2 and an axis y = 0 should be 2 · (0, 0), because if one of the cycles moves (yet in an undefined sense), there are precisely two intersection points which both converge to (0, 0) when the cycles approach the depicted position. (The picture is misleading insofar as the apparently empty intersection of the parabola and the line y = −3 is empty, because only the real solutions of the equations are depicted).\n\nThe first fully satisfactory definition of intersection multiplicities was given by Serre: Let the ambient variety Template:Mvar be smooth (or all local rings regular). Further let Template:Mvar and Template:Mvar be two (irreducible reduced closed) subvarieties, such that their intersection is proper. The construction is local, therefore the varieties may be represented by two ideals Template:Mvar and Template:Mvar in the coordinate ring of Template:Mvar. Let Template:Mvar be an irreducible component of the set-theoretic intersection VW and Template:Mvar its generic point. The multiplicity of Template:Mvar in the intersection product V · W is defined by\n\n$\\mu (Z;V,W):=\\sum _{i=0}^{\\infty }(-1)^{i}{\\text{length}}_{{\\mathcal {O}}_{X,z}}{\\text{Tor}}_{{\\mathcal {O}}_{X,z}}^{i}({\\mathcal {O}}_{X,z}/I,{\\mathcal {O}}_{X,z}/J)$", null, ",\n\nthe alternating sum over the length over the local ring of Template:Mvar in Template:Mvar of torsion groups of the factor rings corresponding to the subvarieties. This expression is sometimes referred to as Serre's Tor-formula.\n\nRemarks:\n\n• The first summand, the length of\n$\\left({\\mathcal {O}}_{X,z}/I\\right)\\otimes _{{\\mathcal {O}}_{X,z}}\\left({\\mathcal {O}}_{X,z}/J\\right)={\\mathcal {O}}_{Z,z}$", null, "is the \"naive\" guess of the multiplicity; however, as Serre shows, it is not sufficient.\n\n### The Chow ring\n\n{{#invoke:main|main}} The Chow ring is the group of algebraic cycles modulo rational equivalence together with the following commutative intersection product:\n\n$V\\cdot W:=\\sum _{i}\\mu (Z_{i};V,W)Z_{i}$", null, "where VW = ∪︀ Zi is the decomposition of the set-theoretic intersection into irreducible components.\n\n### Self-intersection\n\nGiven two subvarieties Template:Mvar and Template:Mvar, one can take their intersection VW, but it is also possible, though more subtle, to define the self-intersection of a single subvariety.\n\nGiven, for instance, a curve Template:Mvar on a surface Template:Mvar, its intersection with itself (as sets) is just itself: CC = C. This is clearly correct, but on the other hand unsatisfactory: given any two distinct curves on a surface (with no component in common), they intersect in some set of points, which for instance one can count, obtaining an intersection number, and we may wish to do the same for a given curve: the analogy is that intersecting distinct curves is like multiplying two numbers: Template:Mvar, while self-intersection is like squaring a single number: x2. Formally, the analogy is stated as a symmetric bilinear form (multiplication) and a quadratic form (squaring).\n\nA geometric solution to this is to intersect the curve Template:Mvar not with itself, but with a slightly pushed off version of itself. In the plane, this just means translating the curve Template:Mvar in some direction, but in general one talks about taking a curve C′ that is linearly equivalent to Template:Mvar, and counting the intersection C · C′, thus obtaining an intersection number, denoted C · C. Note that unlike for distinct curves Template:Mvar and Template:Mvar, the actual points of intersection are not defined, because they depend on a choice of C′, but the “self intersection points of C′′ can be interpreted as Template:Mvar generic points on Template:Mvar, where k = C · C. More properly, the self-intersection points of Template:Mvar is the generic point of Template:Mvar, taken with multiplicity C · C.\n\nAlternatively, one can “solve” (or motivate) this problem algebraically by dualizing, and looking at the class of the [C] ∪ [C] – this both gives a number, and raises the question of a geometric interpretation. Note that passing to cohomology classes is analogous to replacing a curve by a linear system.\n\nNote that self-intersection number can be negative, as the example below illustrates.\n\n#### Examples\n\nConsider a line Template:Mvar in the projective plane P2: it has self-intersection number 1 since all other lines cross it once: one can push Template:Mvar off to L′, and L · L′ = 1 (for any choice) of L′, hence L · L = 1. In terms of intersection forms, we say the plane has one of type x2 (there is only one class of lines, and they all intersect with each other).\n\nNote that on the affine plane, one might push off Template:Mvar to a parallel line, so (thinking geometrically) the number of intersection points depends on the choice of push-off. One says that “the affine plane does not have a good intersection theory”, and intersection theory on non-projective varieties is much more difficult.\n\nA line on a P1 × P1 (which can also be interpreted as the non-singular quadric Template:Mvar in P3) has self-intersection 0, since a line can be moved off itself. (It is a ruled surface.) In terms of intersection forms, we say P1 × P1 has one of type Template:Mvar (which can also be stated x2y2 under a change of basis) – there are two basic classes of lines, which intersect each other in one point (Template:Mvar), but have zero self-intersection (no x2 or y2 terms).\n\n#### Blow-ups\n\nA key example of self-intersection numbers is the exceptional curve of a blow-up, which is a central operation in birational geometry.\n\nGiven an algebraic surface Template:Mvar, blowing up at a point creates a curve Template:Mvar. This curve Template:Mvar is recognisable by its genus, which is 0, and its self-intersection number, which is −1. (This is not obvious.)\n\nNote that as a corollary, P2 and P1 × P1 are minimal surfaces (they are not blow-ups), since they do not have any curves with negative self-intersection.\n\nIn fact, Castelnuovo’s contraction theorem states the converse: every (−1)-curve is the exceptional curve of some blow-up (it can be “blown down”)." ]
[ null, "https://demo.formulasearchengine.com/index.php", null, "https://demo.formulasearchengine.com/index.php", null, "https://demo.formulasearchengine.com/index.php", null, "https://demo.formulasearchengine.com/index.php", null, "https://demo.formulasearchengine.com/index.php", null, "https://demo.formulasearchengine.com/index.php", null ]
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https://www.colorhexa.com/0cf1ff
[ "# #0cf1ff Color Information\n\nIn a RGB color space, hex #0cf1ff is composed of 4.7% red, 94.5% green and 100% blue. Whereas in a CMYK color space, it is composed of 95.3% cyan, 5.5% magenta, 0% yellow and 0% black. It has a hue angle of 183.5 degrees, a saturation of 100% and a lightness of 52.4%. #0cf1ff color hex could be obtained by blending #18ffff with #00e3ff. Closest websafe color is: #00ffff.\n\n• R 5\n• G 95\n• B 100\nRGB color chart\n• C 95\n• M 5\n• Y 0\n• K 0\nCMYK color chart\n\n#0cf1ff color description : Vivid cyan.\n\n# #0cf1ff Color Conversion\n\nThe hexadecimal color #0cf1ff has RGB values of R:12, G:241, B:255 and CMYK values of C:0.95, M:0.05, Y:0, K:0. Its decimal value is 848383.\n\nHex triplet RGB Decimal 0cf1ff `#0cf1ff` 12, 241, 255 `rgb(12,241,255)` 4.7, 94.5, 100 `rgb(4.7%,94.5%,100%)` 95, 5, 0, 0 183.5°, 100, 52.4 `hsl(183.5,100%,52.4%)` 183.5°, 95.3, 100 00ffff `#00ffff`\nCIE-LAB 87.097, -41.695, -20.176 49.651, 70.204, 105.536 0.22, 0.311, 70.204 87.097, 46.32, 205.823 87.097, -65.568, -26.233 83.788, -40.851, -16.028 00001100, 11110001, 11111111\n\n# Color Schemes with #0cf1ff\n\n• #0cf1ff\n``#0cf1ff` `rgb(12,241,255)``\n• #ff1a0c\n``#ff1a0c` `rgb(255,26,12)``\nComplementary Color\n• #0cff94\n``#0cff94` `rgb(12,255,148)``\n• #0cf1ff\n``#0cf1ff` `rgb(12,241,255)``\n• #0c78ff\n``#0c78ff` `rgb(12,120,255)``\nAnalogous Color\n• #ff940c\n``#ff940c` `rgb(255,148,12)``\n• #0cf1ff\n``#0cf1ff` `rgb(12,241,255)``\n• #ff0c78\n``#ff0c78` `rgb(255,12,120)``\nSplit Complementary Color\n• #f1ff0c\n``#f1ff0c` `rgb(241,255,12)``\n• #0cf1ff\n``#0cf1ff` `rgb(12,241,255)``\n• #ff0cf1\n``#ff0cf1` `rgb(255,12,241)``\n• #0cff1a\n``#0cff1a` `rgb(12,255,26)``\n• #0cf1ff\n``#0cf1ff` `rgb(12,241,255)``\n• #ff0cf1\n``#ff0cf1` `rgb(255,12,241)``\n• #ff1a0c\n``#ff1a0c` `rgb(255,26,12)``\n• #00b4bf\n``#00b4bf` `rgb(0,180,191)``\n• #00ccd8\n``#00ccd8` `rgb(0,204,216)``\n• #00e4f2\n``#00e4f2` `rgb(0,228,242)``\n• #0cf1ff\n``#0cf1ff` `rgb(12,241,255)``\n• #26f2ff\n``#26f2ff` `rgb(38,242,255)``\n• #3ff4ff\n``#3ff4ff` `rgb(63,244,255)``\n• #59f5ff\n``#59f5ff` `rgb(89,245,255)``\nMonochromatic Color\n\n# Alternatives to #0cf1ff\n\nBelow, you can see some colors close to #0cf1ff. Having a set of related colors can be useful if you need an inspirational alternative to your original color choice.\n\n• #0cffd0\n``#0cffd0` `rgb(12,255,208)``\n• #0cffe5\n``#0cffe5` `rgb(12,255,229)``\n• #0cfff9\n``#0cfff9` `rgb(12,255,249)``\n• #0cf1ff\n``#0cf1ff` `rgb(12,241,255)``\n• #0cddff\n``#0cddff` `rgb(12,221,255)``\n• #0cc9ff\n``#0cc9ff` `rgb(12,201,255)``\n• #0cb4ff\n``#0cb4ff` `rgb(12,180,255)``\nSimilar Colors\n\n# #0cf1ff Preview\n\nText with hexadecimal color #0cf1ff\n\nThis text has a font color of #0cf1ff.\n\n``<span style=\"color:#0cf1ff;\">Text here</span>``\n#0cf1ff background color\n\nThis paragraph has a background color of #0cf1ff.\n\n``<p style=\"background-color:#0cf1ff;\">Content here</p>``\n#0cf1ff border color\n\nThis element has a border color of #0cf1ff.\n\n``<div style=\"border:1px solid #0cf1ff;\">Content here</div>``\nCSS codes\n``.text {color:#0cf1ff;}``\n``.background {background-color:#0cf1ff;}``\n``.border {border:1px solid #0cf1ff;}``\n\n# Shades and Tints of #0cf1ff\n\nA shade is achieved by adding black to any pure hue, while a tint is created by mixing white to any pure color. In this example, #000b0c is the darkest color, while #f7ffff is the lightest one.\n\n• #000b0c\n``#000b0c` `rgb(0,11,12)``\n• #001e20\n``#001e20` `rgb(0,30,32)``\n• #003033\n``#003033` `rgb(0,48,51)``\n• #004347\n``#004347` `rgb(0,67,71)``\n• #00555a\n``#00555a` `rgb(0,85,90)``\n• #00686e\n``#00686e` `rgb(0,104,110)``\n• #007a82\n``#007a82` `rgb(0,122,130)``\n• #008d95\n``#008d95` `rgb(0,141,149)``\n• #009fa9\n``#009fa9` `rgb(0,159,169)``\n• #00b2bd\n``#00b2bd` `rgb(0,178,189)``\n• #00c4d0\n``#00c4d0` `rgb(0,196,208)``\n• #00d7e4\n``#00d7e4` `rgb(0,215,228)``\n• #00e9f7\n``#00e9f7` `rgb(0,233,247)``\n• #0cf1ff\n``#0cf1ff` `rgb(12,241,255)``\n• #20f2ff\n``#20f2ff` `rgb(32,242,255)``\n• #33f3ff\n``#33f3ff` `rgb(51,243,255)``\n• #47f4ff\n``#47f4ff` `rgb(71,244,255)``\n• #5af6ff\n``#5af6ff` `rgb(90,246,255)``\n• #6ef7ff\n``#6ef7ff` `rgb(110,247,255)``\n• #82f8ff\n``#82f8ff` `rgb(130,248,255)``\n• #95f9ff\n``#95f9ff` `rgb(149,249,255)``\n• #a9faff\n``#a9faff` `rgb(169,250,255)``\n• #bdfbff\n``#bdfbff` `rgb(189,251,255)``\n• #d0fcff\n``#d0fcff` `rgb(208,252,255)``\n• #e4fdff\n``#e4fdff` `rgb(228,253,255)``\n• #f7ffff\n``#f7ffff` `rgb(247,255,255)``\nTint Color Variation\n\n# Tones of #0cf1ff\n\nA tone is produced by adding gray to any pure hue. In this case, #7c8e8f is the less saturated color, while #0cf1ff is the most saturated one.\n\n• #7c8e8f\n``#7c8e8f` `rgb(124,142,143)``\n• #739698\n``#739698` `rgb(115,150,152)``\n• #699ea2\n``#699ea2` `rgb(105,158,162)``\n• #60a7ab\n``#60a7ab` `rgb(96,167,171)``\n• #57afb4\n``#57afb4` `rgb(87,175,180)``\n• #4db7be\n``#4db7be` `rgb(77,183,190)``\n• #44bfc7\n``#44bfc7` `rgb(68,191,199)``\n• #3bc8d0\n``#3bc8d0` `rgb(59,200,208)``\n• #31d0da\n``#31d0da` `rgb(49,208,218)``\n• #28d8e3\n``#28d8e3` `rgb(40,216,227)``\n• #1fe0ec\n``#1fe0ec` `rgb(31,224,236)``\n• #15e9f6\n``#15e9f6` `rgb(21,233,246)``\n• #0cf1ff\n``#0cf1ff` `rgb(12,241,255)``\nTone Color Variation\n\n# Color Blindness Simulator\n\nBelow, you can see how #0cf1ff is perceived by people affected by a color vision deficiency. This can be useful if you need to ensure your color combinations are accessible to color-blind users.\n\nMonochromacy\n• Achromatopsia 0.005% of the population\n• Atypical Achromatopsia 0.001% of the population\nDichromacy\n• Protanopia 1% of men\n• Deuteranopia 1% of men\n• Tritanopia 0.001% of the population\nTrichromacy\n• Protanomaly 1% of men, 0.01% of women\n• Deuteranomaly 6% of men, 0.4% of women\n• Tritanomaly 0.01% of the population" ]
[ null ]
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https://www.vadepares.cat/geometric-series-worksheet
[ "", null, "Home » coloring page » Geometric Series Worksheet\n\n# Geometric Series Worksheet\n\nUploaded by under coloring page [2 views ]\n\n1) a 1 = −3, r = 4 2) a 1 = 4, r = − 3 4 3) a 1 Arithmetic and geometric series worksheet mcr3u jensen general formula for an arithmetic series.", null, "Arithmetic and Geometric Sequences and Series Games\n\n### Arithmetic and geometric series worksheet kuta.", null, "Geometric series worksheet. Some of the worksheets displayed are arithmetic and geometric sequences work comparing arithmetic and geometric sequences arithmetic sequences date period concept 16 arithmetic geometric sequences algebra 2 geometric. Sn = a(1 rn) 1 r example 5 : Khan academy is a 501(c)(3) nonprofit organization.\n\nOut of tips on conversation writing, to cooking publication outlines, as well as to identifying the kind of phrases to use for the. A) 1 1 + 1 1 + 1 1 + b) 1 2 + 3 10 + 9 50 + 27 250 + c) 1 p 2 + 1 2 + 1 p 8 + 1 4 + d) x+ x3 + x9 + x27 + e) :5 + :5 + :5 + :5 + f) x 1 + x2 2 + x3 3 + x4 4 + 2. So an example of a geometric series is 1+ 1 10 + 1 100 + 1 1000 + we can take the sum of the rst n terms of a geometric series and this is denoted by sn:\n\nArithmetic and geometric sequence worksheet worksheet fun worksheet 51 geometric series ks ia2 arithmetic and geometric means kuta software arithmetic and geometric means kuta software kuta software infinite algebra 2 finite geometric series arithmetic sequences. When writing the formula the only thing you fill in is the 1st term and either d or r. This digital resource can be used as an quiz, hw, or paperless worksheet and is designed with google forms™ which means it is no prep for you.\n\nTopics on the quiz include the sum of even integers and the formula for finding an. Given the rst two terms of a geometric progression as 2 and 4, what Knowledge of relevant formulae is a prerequisite to evaluate.\n\n\\( 5 + 25 + 125 + 625… \\) common ratio. In contrast, the power series written as a 0 + a 1 r + a 2 r 2 + a 3 r 3 +. The first term of a geometric series.\n\nFor example, each term in this series is a power of 1/2. 3 + 6 + 12 +. Arithmetic and geometric series worksheet mcr3u jensen general formula for an arithmetic series.\n\nWhich of the following are geometric series? 1) 2, 12 , 72 , 432 2) −1, 5, −25 , 125 Applications of sequences and series worksheet.\n\nEvery coefficient in the geometric series is the same. For each in nite geometric. General formula for a geometric series.\n\nBecause we should supply all you need a single true and also trustworthy supply, many of us found handy info on many themes and topics. Free (33) srwhitehouse a level maths c2: In fact with geometric sequences especially when we re dealing with growth or decay like with money we ll see that they equations look a lot like some of the exponential equations we worked with here in the exponential functions section.\n\nSample our free worksheets and start off your. A series is the sum of the terms of a sequence. An online calculator to calculate the sum of the terms in an arithmetic sequence.\n\nA geometric series is a geometric progression with plus signs between the terms instead of commas. A geometric series is a series or summation that sums the terms of a geometric sequence. In expanded form has coefficients a i that can vary from term to term.\n\nFind the sum of the geometric series. 1 13 15 17 19 21 23 108. Students should have the sequence right before they start the work.\n\nThere are 24 questions which will be grade You can provide visual aids, which can help your students to learn the first step, the second step, and the third step. Some of the worksheets for this concept are finite geometric series, 9 11 sequences word, geometric sequences and series, geometric and arithmetic series word problems, , geometry word problems no problem, arithmetic and geometric series work 1, arithmetic sequences series work.\n\nUse this quiz and worksheet to practice with arithmetic and a geometric series. Geometric sequence and series worksheet along with instructional contents. Whenever there is a constant ratio from one term to the next, the series is called geometric.\n\nThe constant which is multiplied to the preceding term is known as the common ratio,\\( r \\). Some of the worksheets displayed are geometric sequences date period, concept 16 arithmetic geometric sequences, geometric sequence 9nkkzr, geometric sequence, geometric sequence and series work, geometric sequences work, work 3 6 arithmetic and geometric progressions, finite geometric series. For the ones which are geometric series, determine whether they converge.\n\nWhat makes the series geometric is that each term is a power of a constant base. Some of the worksheets for this concept are finite geometric series, infinite geometric series, chapter chapter standardized test, geometric series in financial mathematics, kuta software infinite algebra 2 geometric series, infinite geometric series, geometric sequences and series, work 3 6 arithmetic and. Once you find your worksheet click on pop out.\n\nThe geometric series a + ar + ar 2 + ar 3 +. Apart from the stuff given above, if you want to know more about geometric series worksheet , please click hereapart from the stuff given in this section, if you need any other stuff in math, please use our google custom search here. Found worksheet you are looking for?\n\nA series whose terms are in geometric progression is called geometric series. Geometric sequence and series worksheet. After having gone through the stuff, we hope that the students would have understood geometric series worksheet .\n\nAn array of topics, like evaluating the sum of the geometric series, determining the first term, common ratio and number of terms, exercises on summation notation are included. Some of the worksheets displayed are geometric sequences date period, finite geometric series, geometric sequences work, concept 16 arithmetic geometric sequences, geometric sequences 1314, arithmetic and geometric series work 1, work 3 6 arithmetic and geometric progressions, sequences work 1. Access this finite geometric series worksheets tenaciously prepared for high school students.\n\nAbout this quiz & worksheet. Before referring to geometric sequences and series worksheet answers, you should be aware that knowledge is the answer to a much better the day after tomorrow, in addition to understanding doesn’t just avoid the moment the education bell rings.which remaining mentioned, we all provide you with a number of easy however helpful reports and themes made suitable for almost any helpful purpose. There are methods and formulas we can use to find the value of a geometric series.\n\nIs written in expanded form. In other words, the geometric series is a special case of the power series. O t pmkaadwet 0wmihtahr limngfji sn ciftie o pa el gggehb8rka 6 k2i.\n\n1 a n 40 5n 2 a n 176 200n determine if the sequence is arithmetic. Free (51) srwhitehouse a level maths: \\( 5, 25, 125, 625… \\) geometric series:\n\nFree algebra 2 worksheets created with infinite algebra 2. Geometric sequences and series worksheet answers along with best arithmetic and geometric sequences worksheet awesome 113. Assessing your precalculus students' skills working with geometric sequences and series in has never been easier.\n\nGeometric sequence worksheets are prepared for determining the geometric sequence, finding first term and common ratio, finding the n th term of a geometric sequence, finding next three terms of the sequence and much more. Let a, ar, ar 2,. A geometric progression is a series or sequence where each term is found by multiplying the previous term by a constant:\n\nEquivalently, each term is half of its predecessor.", null, "Geometric Sequences and Series Practice Sequence and", null, "Arithmetic and Geometric Sequences Worksheet Beautiful", null, "Arithmetic & Geometric Sequences Arithmetic, Algebra", null, "Pin on Printable Blank Worksheet Template", null, "50 Geometric Sequences Worksheet Answers in 2020", null, "50 Geometric Sequences Worksheet Answers in 2020", null, "Arithmetic and Geometric Sequences Worksheet 50 Geometric", null, "Pin on Printable Education Worksheet Templates", null, "Geometric Sequences Worksheet Answers Fresh Geometric", null, "Wizer.Me blended worksheet \"Arithmetic & Geometric", null, "50 Arithmetic Sequences and Series Worksheet in 2020", null, "Arithmetic Sequences Maze Worksheet Algebra activities", null, "Arithmetic Sequence Worksheet Answers Math Test Arithmetic", null, "50 Arithmetic Sequences and Series Worksheet Chessmuseum", null, "Geometric sequences worksheet Math worksheet, Math", null, "Geometric Sequences and Series RIDDLE WORKSHEET Sequence", null, "50 Arithmetic Sequences and Series Worksheet in 2020", null, "Pin on Printable Blank Worksheet Template", null, "Geometric Sequences Worksheet Answers Luxury" ]
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https://metanumbers.com/554368
[ "## 554368\n\n554,368 (five hundred fifty-four thousand three hundred sixty-eight) is an even six-digits composite number following 554367 and preceding 554369. In scientific notation, it is written as 5.54368 × 105. The sum of its digits is 31. It has a total of 9 prime factors and 32 positive divisors. There are 268,800 positive integers (up to 554368) that are relatively prime to 554368.\n\n## Basic properties\n\n• Is Prime? No\n• Number parity Even\n• Number length 6\n• Sum of Digits 31\n• Digital Root 4\n\n## Name\n\nShort name 554 thousand 368 five hundred fifty-four thousand three hundred sixty-eight\n\n## Notation\n\nScientific notation 5.54368 × 105 554.368 × 103\n\n## Prime Factorization of 554368\n\nPrime Factorization 27 × 61 × 71\n\nComposite number\nDistinct Factors Total Factors Radical ω(n) 3 Total number of distinct prime factors Ω(n) 9 Total number of prime factors rad(n) 8662 Product of the distinct prime numbers λ(n) -1 Returns the parity of Ω(n), such that λ(n) = (-1)Ω(n) μ(n) 0 Returns: 1, if n has an even number of prime factors (and is square free) −1, if n has an odd number of prime factors (and is square free) 0, if n has a squared prime factor Λ(n) 0 Returns log(p) if n is a power pk of any prime p (for any k >= 1), else returns 0\n\nThe prime factorization of 554,368 is 27 × 61 × 71. Since it has a total of 9 prime factors, 554,368 is a composite number.\n\n## Divisors of 554368\n\n32 divisors\n\n Even divisors 28 4 2 2\nTotal Divisors Sum of Divisors Aliquot Sum τ(n) 32 Total number of the positive divisors of n σ(n) 1.13832e+06 Sum of all the positive divisors of n s(n) 583952 Sum of the proper positive divisors of n A(n) 35572.5 Returns the sum of divisors (σ(n)) divided by the total number of divisors (τ(n)) G(n) 744.559 Returns the nth root of the product of n divisors H(n) 15.5842 Returns the total number of divisors (τ(n)) divided by the sum of the reciprocal of each divisors\n\nThe number 554,368 can be divided by 32 positive divisors (out of which 28 are even, and 4 are odd). The sum of these divisors (counting 554,368) is 1,138,320, the average is 3,557,2.5.\n\n## Other Arithmetic Functions (n = 554368)\n\n1 φ(n) n\nEuler Totient Carmichael Lambda Prime Pi φ(n) 268800 Total number of positive integers not greater than n that are coprime to n λ(n) 6720 Smallest positive number such that aλ(n) ≡ 1 (mod n) for all a coprime to n π(n) ≈ 45542 Total number of primes less than or equal to n r2(n) 0 The number of ways n can be represented as the sum of 2 squares\n\nThere are 268,800 positive integers (less than 554,368) that are coprime with 554,368. And there are approximately 45,542 prime numbers less than or equal to 554,368.\n\n## Divisibility of 554368\n\n m n mod m 2 3 4 5 6 7 8 9 0 1 0 3 4 3 0 4\n\nThe number 554,368 is divisible by 2, 4 and 8.\n\n• Refactorable\n• Abundant\n\n• Polite\n• Practical\n\n• Frugal\n\n## Base conversion (554368)\n\nBase System Value\n2 Binary 10000111010110000000\n3 Ternary 1001011110011\n4 Quaternary 2013112000\n5 Quinary 120214433\n6 Senary 15514304\n8 Octal 2072600\n10 Decimal 554368\n12 Duodecimal 228994\n16 Hexadecimal 87580\n20 Vigesimal 395i8\n36 Base36 bvr4\n\n## Basic calculations (n = 554368)\n\n### Multiplication\n\nn×i\n n×2 1108736 1663104 2217472 2771840\n\n### Division\n\nni\n n⁄2 277184 184789 138592 110874\n\n### Exponentiation\n\nni\n n2 307323879424 170370524388524032 94447966864217290571776 52358930494582410939694317568\n\n### Nth Root\n\ni√n\n 2√n 744.559 82.1485 27.2866 14.0851\n\n## 554368 as geometric shapes\n\n### Circle\n\nRadius = n\n Diameter 1.10874e+06 3.4832e+06 9.65486e+11\n\n### Sphere\n\nRadius = n\n Volume 7.13646e+17 3.86195e+12 3.4832e+06\n\n### Square\n\nLength = n\n Perimeter 2.21747e+06 3.07324e+11 783995\n\n### Cube\n\nLength = n\n Surface area 1.84394e+12 1.70371e+17 960194\n\n### Equilateral Triangle\n\nLength = n\n Perimeter 1.6631e+06 1.33075e+11 480097\n\n### Triangular Pyramid\n\nLength = n\n Surface area 5.32301e+11 2.00784e+16 452640\n\n## Cryptographic Hash Functions\n\nmd5 48f494a795ebd2f6a1a187cfb6615f45 c490f08e082190553cbf2bed423e5f761e84c149 cfe37a8afdeb15a1e2f4e7077cf494f84d8378770102da6aba40e38df683fde3 36e9c296ef58f4fba46b887de26ea43f84dd196b4cc833498a2d5ee430ec730deeef4d3a315334357c6ac57df63fb6a15bdef5a55a000580557310194ecbe691 c651a727a67e26f5d86389116f2365b4eadf8828" ]
[ null ]
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https://hal.inria.fr/inria-00585561
[ "# Circuits in graphs through a prescribed set of ordered vertices\n\n1 MASCOTTE - Algorithms, simulation, combinatorics and optimization for telecommunications\nCRISAM - Inria Sophia Antipolis - Méditerranée , Laboratoire I3S - COMRED - COMmunications, Réseaux, systèmes Embarqués et Distribués\n2 ALGCO - Algorithmes, Graphes et Combinatoire\nLIRMM - Laboratoire d'Informatique de Robotique et de Microélectronique de Montpellier\nAbstract : A circuit in a simple undirected graph G=(V,E) is a sequence of vertices {v_1,v_2,...,v_{k+1}} such that v_1=v_{k+1} and {v_i,v_{i+1}} in E for i=1,...,k. A circuit C is said to be edge-simple if no edge of G is used twice in C. In this article we study the following problem: which is the largest integer k such that, given any subset of k ordered vertices of a graph G, there exists an edge-simple circuit visiting the k vertices in the prescribed order? We first study the case when G has maximum degree at most 3, establishing the value of k for several subcases, such as when G is planar or 3-vertex-connected. Our main result is that k=10 in infinite square grids. To prove this, we introduce a methodology based on the notion of core graph, in order to reduce the number of possible vertex configurations, and then we test each one of the resulting configurations with an Integer Linear Program (ILP) solver.\nMots-clés :\nDocument type :\nJournal articles\n\nCited literature [19 references]\n\nhttps://hal.inria.fr/inria-00585561\nContributor : David Coudert <>\nSubmitted on : Wednesday, April 13, 2011 - 12:06:08 PM\nLast modification on : Wednesday, October 14, 2020 - 4:23:43 AM\nLong-term archiving on: : Saturday, December 3, 2016 - 5:35:05 PM\n\n### File\n\njoin-final-noformat.pdf\nFiles produced by the author(s)\n\n### Citation\n\nDavid Coudert, Frédéric Giroire, Ignasi Sau Valls. Circuits in graphs through a prescribed set of ordered vertices. Journal of Interconnection Networks, World Scientific Publishing, 2011, 11 (3-4), pp.121-141. ⟨10.1142/S0219265910002763⟩. ⟨inria-00585561⟩\n\nRecord views" ]
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http://cafe-am-hafen.de/inductance-snow/inductance-snow.html
[ "# Formula for the Inductance of a Helix, from Paper by Chester Snow, Coded in Lua\n\nI spotted the very interesting, classical paper,\n\nThe original paper is available online at the NIST Virtual Library. This publication for the first time takes care also of the real, helical structure of a coil inductance, providing a higher precision than the former methods for inductance calculation, e. g., by approximating the real coils through current sheets with correction values, or by summation of self and mutual inductances of wire rings and circular filaments. Calculation examples are given in the paper for wires of round and rectangular section.\n\nThe formulas in the paper are rather complex, and their lengthy derivation intractable (for me). For sure they were very tedious to compute at their time of publishing. And the paper tries by judicious neglect of terms of higher order to keep the complexity manageable. Back then the word \"computer\" was already known, but it meant staff teams who were trained to calculate formulas mostly by hand. Integrals (e. g., function B1(k) in the paper) were solved by a mechanical graphing device named integraph, an ingenious precision machine long since forgotten.\n\nThe beauty of the formulas by Chester Snow is in the possibility to arrive at a very precise result for the inductance of a wire coil, considering its helicity and wire diameters. Today modern PCs can calculate complex formulas like the ones by Snow easily, and there is no real need anymore to use any of the imprecise \"handbook formulas\" for coil inductance, which have been published in the meantime.\n\n## 2016: 90th Anniversary of Chester Snow's Seminal Paper\n\nThe date 25 February 2016 is the 90th anniversary of the submission date from Snow's paper. To make his inductance calculation formulas easily usable in today's practice, and to have them accessible for comparison with new developments, i coded them in the language Lua. The program code tries to follow the formulas exactly (while fixing obvious typos), and to reference the original formulas wherever possible. The integraph method for calculation of function `B1(k)` is replaced by Gauss-Legendre integration, using constants calculated by this Lua program. In addition i have tried to recalculate a few constants to a higher precision than given in the paper. But comparison between the original and the new values shows, that in the paper their precision was already chosen to be just high enough not to influence the inductance result notably, and short enough for efficient calculation by hand.\n\nHere now is the fresh Lua program, written as a library:\n\nTo see how it works, just \"`require()`\" it from a Lua prompt and run the `test()` program, or remove the comment dashes in the line \"`-- t.test()`\", and run the code directly by the Lua interpreter from the command line.\n\nI would like to recommend also the webpage Numerical Methods for Inductance Calculation by Robert Weaver, who discusses Snow's helical inductance calculation method from the above paper in detail, and who also provides a fresh calculation method, thriving for utmost precision.\n\nHave fun!\n\nThis page first put online 22 February 2016." ]
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https://de.scribd.com/document/362204349/Assignment-2-Estimation-Exercises-1-2-and-5
[ "You are on page 1of 8\n\n# Assignment #2: Estimation\n\nDue on ?, 2017\n\n## Exercises 7.1, 7.3, 7.7, 7.69 of Newbold (2015)\n\nLEARNING OUTCOMES\n\nAfter completing the second part of this chapter, students will be able\nto:\n\n## (LO 1) explain the desired properties of estimators\n\n(LO 2) explain the difference between point estimates and confidence\nintervals for estimating a population parameter\n(LO 3) construct and interpret confidence interval estimates for the mean\n(variance known and unknown)\n(LO 4) describe the characteristics of Students t distribution.\n(LO 5) use the Students t distribution to solve business problems.\n(LO 6) construct and interpret confidence interval estimates for population\nproportion (large samples)\n(LO 7) construct confidence interval estimates for the variance of a normal\ndistribution\n(LO 8) determine the required sample size to obtain a desired margin of\nerror for large populations\n\nConceptual Questions\n\n1. Explain why the critical value for a given confidence level when the population variance is not\nknown is always greater than the critical value for the same confidence level when the\npopulation variance is known.\n2. When we need to estimate the population mean, and the population standard deviation is\nunknown, we are hit with a double whammy when it comes to the margin of error. Explain\nwhat the double whammy is and why it occurs. (Hint: Consider the sources of variation in\nthe margin of error.)\n3. An insurance company in Iowa recently conducted a survey of its automobile policy\ncustomers to estimate the mean miles these customers commute to work each day. The result\nbased on a random sample of 300 policyholders indicated the population mean was between\n3.5 and 6.7 miles. This interval estimate was constructed using 95% confidence.\n\n## Probability &Business Statistics II (Fall 2017-18) Page 1/8\n\nAfter receiving this result, one of the managers was overheard telling a colleague that 95% of\nall customers commute between 3.5 and 6.7 miles to work each day. How would you respond\nto this statement? Is it correct? Why or why not? Discuss.\n4. Examine the equation for the margin of error when estimating a population mean\n\ne z / 2\nn\n\nIndicate the effect on the margin of error resulting from an increase in each of the following\nitems:\na. confidence level\nb. z-value\nc. standard deviation\nd. sample size\ne. standard error\n5. What are the advantages of using interval estimation rather than point estimation?\n6. Draw a sampling distribution of an unbiased estimator.\n7. Draw a sampling distribution of a biased estimator.\n8. Draw a diagram that shows the sampling distribution representing two unbiased estimators,\none of which is relatively efficient.\n9. The following statement was made by an expert in statistics. If a sample of size 9 is chosen\nfrom a normal distribution having mean , then we can be 95 percent certain that will lie\n\nX 1.96.s / n X\nwithin where is the sample mean and s is the sample standard deviation.\n10. What is the difference between the student and standard normal distribution.\n11. Consider the t-distribution with 10 degrees of freedom, find a so that P (t < a) = 0.95.\n12. Consider the t-distribution with 15 degrees of freedom, find b so that P (t > b) = 0.025.\n13. Consider the t-distribution with 10 degrees of freedom, use the t table to approximate P(t > 3).\n14.A confidence interval was used to estimate the proportion of statistics\nstudents who are female. A random sample of 72 statistics students\ngenerated the following confidence interval: (.438, .642). Using the\ninformation above, what sample size would be necessary if we wanted to\nestimate the true proportion to within 3% using 99% reliability?\nMCQs\n\n## 1. Which statement is not true about confidence intervals?\n\n(a) A confidence interval is an interval of values computed from\nsample data that is likely to include the true population value.\n(b) An approximate formula for a 95% confidence interval is sample\nestimate margin of error.\n(c) A confidence interval between 20% and 40% means that the\npopulation proportion lies between 20% and 40%.\n\n## Probability &Business Statistics II (Fall 2017-18) Page 2/8\n\n(d) A 99% confidence interval procedure has a higher probability of\nproducing intervals that will include the population parameter than a\n95% confidence interval procedure.\n2. Which statement is not true about the 95% confidence level?\n(a) Confidence intervals computed by using the same procedure will\ninclude the true population value for 95% of all possible random\nsamples taken from the population.\n(b) The procedure that is used to determine the confidence interval\nwill provide an interval that includes the population parameter with\nprobability of 0.95.\n(c) The probability that the true value of the population parameter\nfalls between the bounds of an already computed confidence\ninterval is roughly 95%.\n(d) If we consider all possible randomly selected samples of the\nsame size from a population, the 95% is the percentage of those\nsamples for which the confidence interval includes the population\nparameter.\n3. A 90% confidence interval for the mean percentage of airline reservations\nbeing canceled on the day of the flight is (1.1%, 3.2%). What is the point\nestimator of the mean percentage of reservations that are canceled on the\nday of the flight?\n(a)1.05%\n(b) 2.15%\n(c) 2.1%\n(d) 1.60%\n4. The real estate industry claims that it is the best and most effective\nsystem to market residential real estate. A survey of randomly selected\nhome sellers in Illinois found that a 95% confidence interval for the\nproportion of homes that are sold by a real estate agent is 69% to 81%.\nInterpret the interval in this context.\n(a) In 95% of the years, between 69% and 81% of homes in Illinois\nare sold by a real estate agent.\n(b) 95% of all random samples of home sellers in Illinois will show\nthat between 69% and 81% of homes are sold by a real estate\nagent.\n(c) If you sell a home in Illinois, you have an 75% 6% chance of\nusing a real estate agent.\n(d) We are 95% confident that between 69% and 81% of homes in\nthis survey are sold by a real estate agent.\n(e) We are 95% confident, based on this sample, that between 69%\nand 81% of all homes in Illinois are sold by a real estate agent.\n\n## Probability &Business Statistics II (Fall 2017-18) Page 3/8\n\n5. The real estate industry claims that it is the best and most effective\nsystem to market residential real estate. A survey of randomly selected\nhome sellers in Illinois found that a 99% confidence interval for the\nproportion of homes that are sold by a real estate agent is 70% to 80%.\nExplain what \"99% confidence\" means in this context.\n(a) In 99% of the years, between 70% and 80% of homes in Illinois\nare sold by a real estate agent.\n(b) About 99% of all random samples of home sellers in Illinois will\nproduce a confidence interval that contains the true proportion of\nhomes sold by a real estate agent.\n(c) There is a 99% chance that the true proportion of home sellers in\nIllinois who sell their home with a real estate agent is between 70%\nand 80%.\n(d) 99% of home sellers in Illinois will sell their home with a real\nestate agent between 70% and 80% of the time.\n(e) About 99% of all random samples of home sellers in Illinois will\nfind that between 70% and 80% of homes are sold by a real estate\nagent.\n6. Sales of a new line of athletic footwear are crucial to the success of a\ncompany. The company wishes to estimate the average weekly sales of the\nnew footwear to within \\$300 with 90% reliability. The initial sales indicate\nthat the standard deviation of the weekly sales figures is approximately\n\\$1100. How many weeks of data must be sampled for the company to get\nthe information it desires?\n(a) 7 weeks\n(b) 37 weeks\n(c) 23 weeks\n(d) 10,915 weeks\n7. An educator wanted to look at the study habits of university students. As\npart of the research, data was collected for three variables - the amount of\ntime (in hours per week) spent studying, the amount of time (in hours per\nweek) spent playing video games and the GPA - for a sample of 20 male\nuniversity students. As part of the research, a 95% confidence interval for\nthe average GPA of all male university students was calculated to be:\n(2.95, 3.10). Which of the following statements is true?\n(a) In construction of the confidence interval, a t-value with 20\ndegrees of freedom was used.\n(b) In construction of the confidence interval, a z-value was used.\n(c) In construction of the confidence interval, a t-value with 19\ndegrees of freedom was used.\n\n## Probability &Business Statistics II (Fall 2017-18) Page 4/8\n\n(d) In construction of the confidence interval, a z-value with 20\ndegrees of freedom was used.\n8. Find the value of t0 such that the following statement is true: P(-t 0 t t0)\n= .95 where the number of degrees of freedom df is equal to 15\n(a) 2.602\n(b) 2.947\n(c) 2.131\n(d) 1.753\n9. A confidence interval was used to estimate the proportion of statistics\nstudents who are female. A random sample of 72 statistics students\ngenerated the following confidence interval: (.438, .642). Using the\ninformation above, what sample size would be necessary if we wanted to\nestimate the true proportion to within 3% using 99% reliability?\n(a) 1916\n(b) 1831\n(c) 1769\n(d) 1842\n10.What is the confidence level of the following confidence interval for ?\n\nx 2.33 / n\n\n(a) 233%\n(b) 67%\n(c) 98%\n(d) 78%\nExercise #1\n\n## The personnel director of a large corporation wishes to study absenteeism among\n\nclerical workers at the corporations central office during the year. A random\nsample of 25 clerical workers reveals the following:\n\nX 9.7 s 4.0\nAbsenteeism: days, days.\n12 clerical workers were absent more than 10 days.\n\n## 2. Construct a 95% confidence interval estimate for the mean number of\n\nabsences for clerical workers during the year. Interpret. (Hint: Assume that\nthe population is normally distributed)\n\n## 3. Construct a 95% confidence interval estimate for the population proportion\n\nof clerical workers absent more than 10 days during the year. Interpret.\n4. Suppose that the personnel director also wishes to take a survey in a\nbranch office. What sample size is needed to have 95% confidence in\nestimating the population mean absenteeism to within 1.5 days if the\n\n## Probability &Business Statistics II (Fall 2017-18) Page 5/8\n\npopulation standard deviation is estimated to be 4.5 days? ((Hint: Assume\nthat the population is normally distributed)\n\nExercise #2\n\nAccording to Gallups poll on consumer behavior, 36% of Americans say they will\nconsider only cars manufactured by an American company when purchasing a\nnew car. (Data extracted from The Gallup Poll, www.gallup.com, March 31, 2010.)\nIf you select a random sample of 200 Americans,\n\n1. What are the population, the variable of interest and the parameter(s)?\n2. What is the sample statistic(s)?\n3. What is the distribution of the sample statistic defined in question 2. (Hint.\nCheck the conditions to obtain its distribution)\n4. What is the probability that the sample will have between 30% and 40%\nwho say they will consider only cars manufactured by an American\ncompany when purchasing a new car?\n\nExercise #3\n\n## The Human Relations Department of Electronics, Inc., would like to include a\n\ndental plan as part of the benefits package. The question is: How much does a\ntypical employee and his or her family spend per year on dental expenses? A\nsample of 45 employees reveals the mean amount spent last year was \\$1,820,\nwith a standard deviation of \\$660.\n\n1. What are the population, the variable of interest and the parameter(s)?\n2. Construct a 95 percent confidence interval for the population mean.\n3. The information from part (2) was given to the president of Electronics, Inc.\nHe indicated he could afford \\$1,700 of dental expenses per employee. Is it\n\nExercise #4\n\nA market researcher for a consumer electronics company wants to study the television viewing habits\nof residents of a particular area. A random sample of 40 respondents is selected, and each respondent\nis instructed to keep a detailed record of all television viewing in a particular week. The results are as\nfollows:\n\nX 15.3 s 3.8\nViewing time per week: hours, hours.\n27 respondents watch the evening news on at least three week nights.\n\n## Probability &Business Statistics II (Fall 2017-18) Page 6/8\n\n1. Construct a 95% confidence interval estimate for the mean amount of television watched per\nweek in this area.\n2. Construct a 95% confidence interval estimate for the population proportion who watch the\nevening news on at least three weeknights per week.\n\nSuppose that the market researcher wants to take another survey in a different location. Answer this\nquestion:\n\n3. What sample size is required to be 95% confident of estimating the population mean viewing\ntime to within hours assuming that the population standard deviation is equal to five hours?\n\nExercise #5\n\n## The sponsors of television shows targeted at the childrens market wanted to\n\nknow the amount of time children spend watching television because the types\nand number of programs and commercials are greatly influenced by this\ninformation. As a result, it was decided to survey 100 North American children\nand ask them to keep track of the number of hours of television they watch each\nweek. The results of the survey are\nDescriptive Statistics: Time\n\n## Variable N N* Mean StDev Minimum Maximum\n\nTime 100 0 27,191 8,373 9,500 50,300\n\n1. What are the population, the variable of interest and the objective of the\nsurvey?\nThe television sponsors want an estimate of the amount of television watched by\nthe average North American child. A confidence level of 95% is judged to be\nappropriate.\n2. What are the population parameter(s) and sample statistic(s) considered to\nreach the objective of the survey.\n3. What is the distribution of the random variable that will be used to\nconstruct the confidence interval in question 4. (Hint. State the\ncondition(s) to obtain its distribution)\n4. Construct a 95% confidence interval estimate of the parameter defined in\nquestion 2.\n5. Interpret the confidence interval estimate obtained in question 4.\n6. Use the following histogram to check the assumption stated in question 3.\n\n## Probability &Business Statistics II (Fall 2017-18) Page 7/8\n\n7. Some people erroneously interpret the confidence interval estimate to\n\nIC1 ( ) a, b 1\nmean, , that there is probability that the population\nmean lies between a and b. Explain.\n\nExercise #6\n\n## A company that produces universal remote controls wanted to determine the\n\nnumber of remote control devices American homes contain. It is assumed that\nthis number is normally distributed. The company hired a statistician to survey\n240 randomly selected homes and determine the number of remote controls. The\nresults are reported in Table 3.\n\n## 1. Define the population, variable of interest and the parameters.\n\n2. Estimate with 95% confidence the average number of remote controls in\nthe United States.\n3. If there are 100 million households, estimate with 95% confidence the total\nnumber of remote controls in the United States.\n\n## Remotes 240 0 9 4,66 2,368\n\nValid N (listwise) 240" ]
[ null ]
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https://math.libretexts.org/Courses/Monroe_Community_College/MTH_098_Elementary_Algebra/4%3A_Graphs/4.6%3A_Find_the_Equation_of_a_Line
[ "# 4.6: Find the Equation of a Line\n\n•", null, "• OpenStax\n• OpenStax\n$$\\newcommand{\\vecs}{\\overset { \\rightharpoonup} {\\mathbf{#1}} }$$ $$\\newcommand{\\vecd}{\\overset{-\\!-\\!\\rightharpoonup}{\\vphantom{a}\\smash {#1}}}$$$$\\newcommand{\\id}{\\mathrm{id}}$$ $$\\newcommand{\\Span}{\\mathrm{span}}$$ $$\\newcommand{\\kernel}{\\mathrm{null}\\,}$$ $$\\newcommand{\\range}{\\mathrm{range}\\,}$$ $$\\newcommand{\\RealPart}{\\mathrm{Re}}$$ $$\\newcommand{\\ImaginaryPart}{\\mathrm{Im}}$$ $$\\newcommand{\\Argument}{\\mathrm{Arg}}$$ $$\\newcommand{\\norm}{\\| #1 \\|}$$ $$\\newcommand{\\inner}{\\langle #1, #2 \\rangle}$$ $$\\newcommand{\\Span}{\\mathrm{span}}$$ $$\\newcommand{\\id}{\\mathrm{id}}$$ $$\\newcommand{\\Span}{\\mathrm{span}}$$ $$\\newcommand{\\kernel}{\\mathrm{null}\\,}$$ $$\\newcommand{\\range}{\\mathrm{range}\\,}$$ $$\\newcommand{\\RealPart}{\\mathrm{Re}}$$ $$\\newcommand{\\ImaginaryPart}{\\mathrm{Im}}$$ $$\\newcommand{\\Argument}{\\mathrm{Arg}}$$ $$\\newcommand{\\norm}{\\| #1 \\|}$$ $$\\newcommand{\\inner}{\\langle #1, #2 \\rangle}$$ $$\\newcommand{\\Span}{\\mathrm{span}}$$$$\\newcommand{\\AA}{\\unicode[.8,0]{x212B}}$$\n\n##### Learning Objectives\n\nBy the end of this section, you will be able to:\n\n• Find an equation of the line given the slope and y-intercept\n• Find an equation of the line given the slope and a point\n• Find an equation of the line given two points\n• Find an equation of a line parallel to a given line\n• Find an equation of a line perpendicular to a given line\n##### Note\n\nBefore you get started, take this readiness quiz.\n\n1. Solve: $$\\frac{2}{3} = \\frac{x}{5}$$.\nIf you missed this problem, review Exercise 2.2.4.\n2. Simplify: $$−\\frac{2}{5}(x−15)$$.\nIf you missed this problem, review Exercise 1.10.34.\n\nHow do online retailers know that ‘you may also like’ a particular item based on something you just ordered? How can economists know how a rise in the minimum wage will affect the unemployment rate? How do medical researchers create drugs to target cancer cells? How can traffic engineers predict the effect on your commuting time of an increase or decrease in gas prices? It’s all mathematics.\n\nYou are at an exciting point in your mathematical journey as the mathematics you are studying has interesting applications in the real world.\n\nThe physical sciences, social sciences, and the business world are full of situations that can be modeled with linear equations relating two variables. Data is collected and graphed. If the data points appear to form a straight line, an equation of that line can be used to predict the value of one variable based on the value of the other variable.\n\nTo create a mathematical model of a linear relation between two variables, we must be able to find the equation of the line. In this section we will look at several ways to write the equation of a line. The specific method we use will be determined by what information we are given.\n\n## Find an Equation of the Line Given the Slope and y-Intercept\n\nWe can easily determine the slope and intercept of a line if the equation was written in slope–intercept form, y=mx+b. Now, we will do the reverse—we will start with the slope and y-intercept and use them to find the equation of the line.\n\n##### Exercise $$\\PageIndex{1}$$\n\nFind an equation of a line with slope −7 and y-intercept (0,−1).\n\nSince we are given the slope and y-intercept of the line, we can substitute the needed values into the slope–intercept form, y=mx+b.\n\n Name the slope.", null, "Name the y-intercept.", null, "Substitute the values into y=mx+b.", null, "", null, "", null, "##### Exercise $$\\PageIndex{2}$$\n\nFind an equation of a line with slope $$\\frac{2}{5}$$ and y-intercept (0,4).\n\n$$y = \\frac{2}{5}x + 4$$\n\n##### Exercise $$\\PageIndex{3}$$\n\nFind an equation of a line with slope −1 and y-intercept (0,−3).\n\n$$y=−x−3$$\n\nSometimes, the slope and intercept need to be determined from the graph.\n\n##### Exercise $$\\PageIndex{4}$$\n\nFind the equation of the line shown.", null, "We need to find the slope and y-intercept of the line from the graph so we can substitute the needed values into the slope–intercept form, y=mx+by=mx+b.\n\nTo find the slope, we choose two points on the graph.\n\nThe y-intercept is (0,−4) and the graph passes through (3,−2).\n\n Find the slope by counting the rise and run.", null, "", null, "Find the y-intercept.", null, "Substitute the values into y=mx+b.", null, "", null, "##### Exercise $$\\PageIndex{5}$$\n\nFind the equation of the line shown in the graph.", null, "$$y=\\frac{3}{5}x+1$$\n\n##### Exercise $$\\PageIndex{6}$$\n\nFind the equation of the line shown in the graph.", null, "$$y=\\frac{4}{3}x−5$$\n\n## Find an Equation of the Line Given the Slope and a Point\n\nFinding an equation of a line using the slope–intercept form of the equation works well when you are given the slope and y-intercept or when you read them off a graph. But what happens when you have another point instead of the y-intercept?\n\nWe are going to use the slope formula to derive another form of an equation of the line. Suppose we have a line that has slope mm and that contains some specific point $$(x_{1}, y_{1})$$ and some other point, which we will just call (x,y). We can write the slope of this line and then change it to a different form.\n\n$$\\begin{array} {lrll}&m &=\\frac{y-y_{1}}{x-x_{1}} \\\\ \\text{Multiply both sides of the equation by }x−x_{1}.&m\\left(x-x_{1}\\right) &=\\left(\\frac{y-y_{1}}{x-x_{1}}\\right)\\left(x-x_{1}\\right) \\\\ \\text{Simplify.}&m\\left(x-x_{1}\\right) &=y-y_{1} \\\\ \\text{Rewrite the equation with the y terms on the left.} &y-y_{1} &=m\\left(x-x_{1}\\right) \\end{array}$$\n\nThis format is called the point–slope form of an equation of a line.\n\n##### POINT–SLOPE FORM OF AN EQUATION OF A LINE\n\nThe point–slope form of an equation of a line with slope mm and containing the point $$(x_{1}, y_{1})$$ is", null, "We can use the point–slope form of an equation to find an equation of a line when we are given the slope and one point. Then we will rewrite the equation in slope–intercept form. Most applications of linear equations use the the slope–intercept form.\n\n##### Exercise $$\\PageIndex{7}$$: Find an Equation of a Line Given the Slope and a Point\n\nFind an equation of a line with slope $$m=\\frac{2}{5}$$ that contains the point (10,3). Write the equation in slope–intercept form.", null, "", null, "", null, "", null, "##### Exercise $$\\PageIndex{8}$$\n\nFind an equation of a line with slope $$m=\\frac{5}{6}$$ and containing the point (6,3).\n\n$$y=\\frac{5}{6}x−2$$\n\n##### Exercise $$\\PageIndex{9}$$\n\nFind an equation of a line with slope $$m=\\frac{2}{3}$$ and containing the point (9,2).\n\n$$y=\\frac{2}{3}x−4$$\n\n##### FIND AN EQUATION OF A LINE GIVEN THE SLOPE AND A POINT.\n1. Identify the slope.\n2. Identify the point.\n3. Substitute the values into the point-slope form, $$y−y_{1}=m(x−x_{1})$$.\n4. Write the equation in slope–intercept form.\n##### Exercise $$\\PageIndex{10}$$\n\nFind an equation of a line with slope $$m=−\\frac{1}{3}$$ that contains the point (6,−4). Write the equation in slope–intercept form.\n\nSince we are given a point and the slope of the line, we can substitute the needed values into the point–slope form, $$y−y_{1}=m(x−x_{1})$$.\n\n Identify the slope.", null, "Identify the point.", null, "Substitute the values into $$y−y_{1}=m(x−x_{1})$$.", null, "", null, "Simplify.", null, "Write in slope–intercept form.", null, "##### Exercise $$\\PageIndex{11}$$\n\nFind an equation of a line with slope $$m=−\\frac{2}{5}$$ and containing the point (10,−5).\n\n$$y=−\\frac{2}{5}x−1$$\n\n##### Exercise $$\\PageIndex{12}$$\n\nFind an equation of a line with slope $$m=−\\frac{3}{4}$$, and containing the point (4,−7).\n\n$$y=−\\frac{3}{4}x−4$$\n\n##### Exercise $$\\PageIndex{13}$$\n\nFind an equation of a horizontal line that contains the point (−1,2). Write the equation in slope–intercept form.\n\nEvery horizontal line has slope 0. We can substitute the slope and points into the point–slope form, $$y−y_{1}=m(x−x_{1})$$.\n\n Identify the slope.", null, "Identify the point.", null, "Substitute the values into $$y−y_{1}=m(x−x_{1})$$.", null, "", null, "Simplify.", null, "", null, "", null, "Write in slope–intercept form. It is in y-form, but could be written y=0x+2.\nDid we end up with the form of a horizontal line, y=a?\n##### Exercise $$\\PageIndex{14}$$\n\nFind an equation of a horizontal line containing the point (−3,8).\n\ny = 8\n\n##### Exercise $$\\PageIndex{15}$$\n\nFind an equation of a horizontal line containing the point (−1,4).\n\ny = 4\n\n## Find an Equation of the Line Given Two Points\n\nWhen real-world data is collected, a linear model can be created from two data points. In the next example we’ll see how to find an equation of a line when just two points are given.\n\nWe have two options so far for finding an equation of a line: slope–intercept or point–slope. Since we will know two points, it will make more sense to use the point–slope form.\n\nBut then we need the slope. Can we find the slope with just two points? Yes. Then, once we have the slope, we can use it and one of the given points to find the equation.\n\n##### Exercise $$\\PageIndex{16}$$: Find an Equation of a Line Given Two Points\n\nFind an equation of a line that contains the points (5,4) and (3,6). Write the equation in slope–intercept form.", null, "", null, "", null, "", null, "Use the point (3,6) and see that you get the same equation.\n\n##### Exercise $$\\PageIndex{17}$$\n\nFind an equation of a line containing the points (3,1) and (5,6).\n\n$$y=\\frac{5}{2}x−\\frac{13}{2}$$\n\n##### Exercise $$\\PageIndex{18}$$\n\nFind an equation of a line containing the points (1,4) and (6,2).\n\n$$y=−\\frac{2}{5}x+\\frac{22}{5}$$\n\n##### FIND AN EQUATION OF A LINE GIVEN TWO POINTS.\n1. Find the slope using the given points.\n2. Choose one point.\n3. Substitute the values into the point-slope form, $$y−y_{1}=m(x−x_{1})$$.\n4. Write the equation in slope–intercept form.\n##### Exercise $$\\PageIndex{19}$$\n\nFind an equation of a line that contains the points (−3,−1) and (2,−2). Write the equation in slope–intercept form.\n\nSince we have two points, we will find an equation of the line using the point–slope form. The first step will be to find the slope.\n\n Find the slope of the line through (−3, −1) and (2, −2).", null, "", null, "", null, "", null, "Choose either point.", null, "Substitute the values into $$y−y_{1}=m(x−x_{1})$$.", null, "", null, "", null, "Write in slope–intercept form.", null, "##### Exercise $$\\PageIndex{20}$$\n\nFind an equation of a line containing the points (−2,−4) and (1,−3).\n\n$$y=\\frac{1}{3}x−\\frac{10}{3}$$\n\n##### Exercise $$\\PageIndex{21}$$\n\nFind an equation of a line containing the points (−4,−3) and (1,−5).\n\n$$y=−\\frac{2}{5}x−\\frac{23}{5}$$\n\n##### Exercise $$\\PageIndex{22}$$\n\nFind an equation of a line that contains the points (−2,4) and (−2,−3). Write the equation in slope–intercept form.\n\nAgain, the first step will be to find the slope.\n\n$$\\begin{array}{lrl} \\text { Find the slope of the line through }(-2,4) \\text { and }(-2,-3) & & &\\\\ &m &=&\\frac{y_{2}-x_{1}}{x_{2}-x_{1}} \\\\ &m &=&\\frac{-3-4}{-2-(-2)} \\\\ &m &= &\\frac{-7}{0} \\\\ \\\\ \\text { The slope is undefined. } & & &\\end{array}$$\n\nThis tells us it is a vertical line. Both of our points have an x-coordinate of −2. So our equation of the line is x=−2. Since there is no yy, we cannot write it in slope–intercept form.\n\nYou may want to sketch a graph using the two given points. Does the graph agree with our conclusion that this is a vertical line?\n\n##### Exercise $$\\PageIndex{23}$$\n\nFind an equation of a line containing the points (5,1) and (5,−4).\n\nx = 5\n\n##### Exercise $$\\PageIndex{24}$$\n\nFind an equation of a line containing the points (−4,4) and (−4,3).\n\nx=−4\n\nWe have seen that we can use either the slope–intercept form or the point–slope form to find an equation of a line. Which form we use will depend on the information we are given. This is summarized in Table $$\\PageIndex{1}$$.\n\n To Write an Equation of a Line If given: Use: Form: Slope and y-intercept slope–intercept y=mx+b Slope and a point point–slope $$y−y_{1}=m(x−x_{1})$$ Two points point–slope $$y−y_{1}=m(x−x_{1})$$\n\n## Find an Equation of a Line Parallel to a Given Line\n\nSuppose we need to find an equation of a line that passes through a specific point and is parallel to a given line. We can use the fact that parallel lines have the same slope. So we will have a point and the slope—just what we need to use the point–slope equation.\n\nFirst let’s look at this graphically.\n\nThe graph shows the graph of y=2x−3. We want to graph a line parallel to this line and passing through the point (−2,1).", null, "Figure $$\\PageIndex{1}$$\n\nWe know that parallel lines have the same slope. So the second line will have the same slope as y=2x−3. That slope is $$m_{\\|} = 2$$. We’ll use the notation $$m_{\\|}$$ to represent the slope of a line parallel to a line with slope m. (Notice that the subscript ∥ looks like two parallel lines.)\n\nThe second line will pass through (−2,1) and have m=2. To graph the line, we start at (−2,1) and count out the rise and run. With m=2 (or $$m=\\frac{2}{1}$$), we count out the rise 2 and the run 1. We draw the line.", null, "Figure $$\\PageIndex{2}$$\n\nDo the lines appear parallel? Does the second line pass through (−2,1)?\n\nNow, let’s see how to do this algebraically.\n\nWe can use either the slope–intercept form or the point–slope form to find an equation of a line. Here we know one point and can find the slope. So we will use the point–slope form.\n\n##### Exercise $$\\PageIndex{25}$$: How to Find an Equation of a Line Parallel to a Given Line\n\nFind an equation of a line parallel to y=2x−3 that contains the point (−2,1). Write the equation in slope–intercept form.", null, "", null, "", null, "", null, "", null, "Does this equation make sense? What is the y-intercept of the line? What is the slope?\n\n##### Exercise $$\\PageIndex{26}$$\n\nFind an equation of a line parallel to the line y=3x+1 that contains the point (4,2). Write the equation in slope–intercept form.\n\ny=3x−10\n\n##### Exercise $$\\PageIndex{27}$$\n\nFind an equation of a line parallel to the line $$y=\\frac{1}{2}x−3$$ that contains the point (6,4).\n\n$$y=\\frac{1}{2}x+1$$\n\n##### FIND AN EQUATION OF A LINE PARALLEL TO A GIVEN LINE.\n1. Find the slope of the given line.\n2. Find the slope of the parallel line.\n3. Identify the point.\n4. Substitute the values into the point–slope form, $$y−y_{1}=m(x−x_{1})$$.\n5. Write the equation in slope–intercept form.\n\n## Find an Equation of a Line Perpendicular to a Given Line\n\nNow, let’s consider perpendicular lines. Suppose we need to find a line passing through a specific point and which is perpendicular to a given line. We can use the fact that perpendicular lines have slopes that are negative reciprocals. We will again use the point–slope equation, like we did with parallel lines.\n\nThe graph shows the graph of y=2x−3. Now, we want to graph a line perpendicular to this line and passing through (−2,1).", null, "Figure $$\\PageIndex{3}$$\n\nWe know that perpendicular lines have slopes that are negative reciprocals. We’ll use the notation $$m_{\\perp}$$ to represent the slope of a line perpendicular to a line with slope m. (Notice that the subscript $$_{\\perp}$$ looks like the right angles made by two perpendicular lines.)\n\n$\\begin{array}{cl}{y=2 x-3} & {\\text { perpendicular line }} \\\\ {m=2} & {m_{\\perp}=-\\frac{1}{2}}\\end{array}$\n\nWe now know the perpendicular line will pass through (−2,1) with $$m_{\\perp}=−\\frac{1}{2}$$.\n\nTo graph the line, we will start at (−2,1) and count out the rise −1 and the run 2. Then we draw the line.", null, "Figure $$\\PageIndex{4}$$\n\nDo the lines appear perpendicular? Does the second line pass through (−2,1)?\n\nNow, let’s see how to do this algebraically. We can use either the slope–intercept form or the point–slope form to find an equation of a line. In this example we know one point, and can find the slope, so we will use the point–slope form.\n\n##### Exercise $$\\PageIndex{28}$$\n\nFind an equation of a line perpendicular to y=2x−3 that contains the point (−2,1). Write the equation in slope–intercept form.", null, "", null, "", null, "", null, "", null, "##### Exercise $$\\PageIndex{29}$$\n\nFind an equation of a line perpendicular to the line y=3x+1 that contains the point (4,2). Write the equation in slope–intercept form.\n\n$$y=−\\frac{1}{3}x+\\frac{10}{3}$$\n\n##### Exercise $$\\PageIndex{30}$$\n\nFind an equation of a line perpendicular to the line $$y=\\frac{1}{2}x−3$$ that contains the point (6,4).\n\ny=−2x+16\n\n##### FIND AN EQUATION OF A LINE PERPENDICULAR TO A GIVEN LINE.\n1. Find the slope of the given line.\n2. Find the slope of the perpendicular line.\n3. Identify the point.\n4. Substitute the values into the point–slope form, $$y−y_{1}=m(x−x_{1})$$.\n5. Write the equation in slope–intercept form.\n##### Exercise $$\\PageIndex{31}$$\n\nFind an equation of a line perpendicular to x=5 that contains the point (3,−2). Write the equation in slope–intercept form.\n\nAgain, since we know one point, the point–slope option seems more promising than the slope–intercept option. We need the slope to use this form, and we know the new line will be perpendicular to x=5. This line is vertical, so its perpendicular will be horizontal. This tells us the $$m_{\\perp}=0$$.\n\n$$\\begin{array}{lrll}{\\text { Identify the point. }} &{(3}&{,}&{-2)}\\\\ {\\text { Identify the slope of the perpendicular line. }} & {m_{\\perp}}&{=}&{0} \\\\ {\\text { Substitute the values into } y-y_{1}=m\\left(x-x_{1}\\right) .} & {y-y_{1}}&{=}&{m\\left(x-x_{1}\\right)} \\\\{} &{y−(−2)}&{=}&{0(x−3)} \\\\{\\text { Simplify. }} & {y+2}&{=}&{0} \\\\ &{y}&{=}&{-2}\\end{array}$$\n\nSketch the graph of both lines. Do they appear to be perpendicular?\n\n##### Exercise $$\\PageIndex{32}$$\n\nFind an equation of a line that is perpendicular to the line x=4 that contains the point (4,−5). Write the equation in slope–intercept form.\n\ny=−5\n\n##### Exercise $$\\PageIndex{33}$$\n\nFind an equation of a line that is perpendicular to the line x=2 that contains the point (2,−1). Write the equation in slope–intercept form.\n\ny=−1\n\nIn Exercise $$\\PageIndex{31}$$, we used the point–slope form to find the equation. We could have looked at this in a different way.\n\nWe want to find a line that is perpendicular to x=5 that contains the point (3,−2). The graph shows us the line x=5 and the point (3,−2).", null, "Figure $$\\PageIndex{5}$$\n\nWe know every line perpendicular to a vertical line is horizontal, so we will sketch the horizontal line through (3,−2).", null, "Figure $$\\PageIndex{6}$$\n\nDo the lines appear perpendicular?\n\nIf we look at a few points on this horizontal line, we notice they all have y-coordinates of −2. So, the equation of the line perpendicular to the vertical line x=5 is y=−2.\n\n##### Exercise $$\\PageIndex{34}$$\n\nFind an equation of a line that is perpendicular to y=−4 that contains the point (−4,2). Write the equation in slope–intercept form.\n\nThe line y=−4 is a horizontal line. Any line perpendicular to it must be vertical, in the form x=a. Since the perpendicular line is vertical and passes through (−4,2), every point on it has an x-coordinate of −4. The equation of the perpendicular line is x=−4. You may want to sketch the lines. Do they appear perpendicular?\n\n##### Exercise $$\\PageIndex{35}$$\n\nFind an equation of a line that is perpendicular to the line y=1 that contains the point (−5,1). Write the equation in slope–intercept form.\n\nx=−5\n\n##### Exercise $$\\PageIndex{36}$$\n\nFind an equation of a line that is perpendicular to the line y=−5 that contains the point (−4,−5).\n\nx=−4\n\n##### Note\n\nAccess this online resource for additional instruction and practice with finding the equation of a line.\n\n## Key Concepts\n\n• To Find an Equation of a Line Given the Slope and a Point\n1. Identify the slope.\n2. Identify the point.\n3. Substitute the values into the point-slope form, $$y−y_{1}=m(x−x_{1})$$.\n4. Write the equation in slope-intercept form.\n• To Find an Equation of a Line Given Two Points\n1. Find the slope using the given points.\n2. Choose one point.\n3. Substitute the values into the point-slope form, $$y−y_{1}=m(x−x_{1})$$.\n4. Write the equation in slope-intercept form.\n• To Write and Equation of a Line\n• If given slope and $$y$$-intercept, use slope–intercept form $$y=mx+b$$.\n• If given slope and a point, use point–slope form $$y−y_{1}=m(x−x_{1})$$.\n• If given two points, use point–slope form $$y−y_{1}=m(x−x_{1})$$.\n• To Find an Equation of a Line Parallel to a Given Line\n1. Find the slope of the given line.\n2. Find the slope of the parallel line.\n3. Identify the point.\n4. Substitute the values into the point-slope form, $$y−y_{1}=m(x−x_{1})$$.\n5. Write the equation in slope-intercept form.\n• To Find an Equation of a Line Perpendicular to a Given Line\n1. Find the slope of the given line.\n2. Find the slope of the perpendicular line.\n3. Identify the point.\n4. Substitute the values into the point-slope form, $$y−y_{1}=m(x−x_{1})$$.\n5. Write the equation in slope-intercept form.\n\n## Glossary\n\npoint–slope form\nThe point–slope form of an equation of a line with slope mm and containing the point $$\\left(x_{1}, y_{1}\\right)$$ is $$y-y_{1}=m\\left(x-x_{1}\\right)$$.\n\nThis page titled 4.6: Find the Equation of a Line is shared under a CC BY license and was authored, remixed, and/or curated by OpenStax." ]
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https://journals.ametsoc.org/mwr/article/138/11/4026/70899/Synoptic-Eddy-Feedbacks-and-Circulation-Regime
[ "## Abstract\n\nA method to incorporate synoptic eddies into the diagnosis of circulation regimes using cluster analysis is illustrated using boreal winter reanalyses of the National Centers of Environmental Prediction (hereafter observations) over the Pacific–North American region. The motivation is to include the configuration of the high-frequency (periods less than 10 days) transients as well as the low-frequency (periods greater than 10 days) flow explicitly into the definition of the regimes.\n\nPrinciple component analysis is applied to the low-frequency 200-hPa height field, and also to the low-frequency “envelope” modulations of the rms of high-frequency meridional velocity at 200 hPa. A maximum covariance analysis of the height and envelope fields, carried out using the appropriate principal components, defines three modes as explaining most of the covariance. This defines the minimum dimensionality of the space in which to apply k-means cluster analysis to the covariance coefficients. Clusters found using this method agree with results of the previous work.\n\nSignificance is assessed by comparing cluster analyses with results from synthetic datasets that have the same spectral amplitudes (but random phases) of seasonal means and, separately, intraseasonal fluctuations as do the original observed time series. This procedure ensures that the synthetic series have similar autocovariance structures to the observations. Building on earlier work, the clusters obtained are newly tested to be highly significant without the need for quasi-stationary prefiltering.\n\n## 1. Introduction\n\nThe analysis of low-frequency atmospheric variability within the framework of circulation regimes has proved quite useful in terms of relating synoptic behavior and extreme surface anomalies to the large-scale flow (Ayrault et al. 1995; Arnott et al. 2004; Cassou et al. 2005; Stan and Straus 2007). Methods that have been used to identify such regimes include cluster analysis (Michelangeli et al. 1995), modeling the probability distribution function in terms of additive Gaussian structures (Smyth et al. 1999), and statistical models based on neural networks (Monahan et al. 2001), although this last method has been questioned (Christiansen 2005).\n\nHowever, the vast majority of observational studies identifying regimes have one feature in common: the analysis is carried out on time series of a single variable (often geopotential height or streamfunction), often filtered to remove synoptic time-scale variability, and expressed in a reduced dimensional representation [usually principal components (PCs)].\n\nTheoretical studies, on the other hand, have shown the importance of baroclinic-scale transient feedback on the large scale as a mechanism to help support circulation regimes (Reinhold and Pierrehumbert 1982; Vautard and Legras 1988). Observational studies have confirmed the feedback between baroclinic waves and low-frequency structures (Nakamura and Wallace 1990; Nakamura et al. 1997).\n\nThe purpose of this article is to introduce a method to identify circulation regimes that includes baroclinic-scale properties. This is accomplished in the context of the work of Straus et al. (2007, hereafter SCM), in which circulation regimes were associated with significant clustering of low-frequency filtered height maps in a state space defined by PCs.\n\nThis paper uses the same reanalysis dataset as in SCM, but incorporates the baroclinic eddy properties to help define a minimum state space dimension. In addition we use an improved method of assessing significance, involving a generalization of the “random phase” approximation of Christiansen (2007).\n\nSection 2 briefly describes the dataset, the preprocessing used, the construction of the envelope function of eddy properties and the maximum covariance analysis. Section 3 describes the cluster analysis and the tests of significance. Brief conclusions are presented in section 4.\n\n## 2. Data and analysis\n\nTime series of once-daily geopotential height and meridional wind at 200 hPa were obtained from the reanalyses of the National Centers for Environmental Prediction (NCEP) for the 152-day period starting 21 November, for each of the 54 winters 1948/49 through 2001/02. An estimated climatological annual cycle was subtracted (as in SCM), so that the resulting series retain both interannual and intraseasonal variability. A digital filter was used to divide the time series of each variable into two components: a “low frequency” component containing all periods longer than about 10 days, and a “high frequency” component (retaining periods shorter than 10 days), representative of baroclinic eddies. The digital filter is similar to that used by, for example, Blackmon (1976), who use a filter that retains periods of 2–6 days. Here the high-frequency category extends up to periods of 10 days, since this filter is meant to capture not only the growth and propagation of baroclinic waves, but their life cycle as well (Simmons and Hoskins 1978). We denote the high-frequency components by primes.\n\nDaily time series of the rms of high-frequency meridional wind were computed as\n\nThe modulation of υrms series on time scales longer than 10 days was obtained by retaining only the low-frequency component. The use of such an “envelope” functions is given (e.g., in Nakamura and Wallace 1990). Since the meridional wind is dominated by the longitudinal derivative of the streamfunction, it emphasizes the higher zonal wavenumbers.\n\nNote that we need to filter twice: once to obtain the high-frequency fields appearing on the right-hand side of Eq. (1) and a second time to low-pass filter the resulting rms velocity. Since each application of the filter results is a loss of 15 days of the record at the beginning and at the end, the high-frequency fields are available for 120 days starting 7 December, but the envelope fields are available only for the 92 days starting 21 December for each winter.\n\nA principal component (or empirical orthogonal function) analysis was carried out on the 92-day time series of low-frequency 200-hPa geopotential height, over the domain 20°–90°N, 150°E–30°W, denoted by the vector Z, and, separately on the low-frequency envelope of υrms, denoted by the vector E. The spectrum of explained variance for Z is fairly red, with the leading three modes explaining 22%, 16%, and 12% of the total space–time variance, respectively, while the leading 10 modes together explain 84%. In contrast, the spectrum of the E modes is somewhat less red, with the leading three modes explaining about 15%, 10%, and 9% of the total variance, and the leading 10 modes explaining 59%.\n\nThe covariance of the envelope field with the low-frequency flow is assessed with the use of maximum covariance analysis (MCA),1 applied to the PCs of Z and E. This analysis identifies linear combinations of PCs that maximize the squared covariance between the two fields. For more details on this and related multivariate analysis techniques, see Bretherton et al. (1992) and Wallace et al. (1992).\n\nThe results of the MCA calculations depend on the truncation (number of PCs) of the input fields. Table 1 gives a number of statistics as a function of PC truncation: the total squared covariance explained, the squared covariance explained by each of the first three modes, the percent of total squared covariance explained, and the percentage of the truncated variance of Z explained by each mode. In addition, the correlations between the leading PCs of Z and E with the corresponding leading MCA time series are shown. The correlations for the first mode are highly significant, indicating that the MCA analysis is not overly dominated by the Z spectrum, which is the redder of the two.2\n\nTable 1.\n\nSummary of maximum covariance calculations between low-frequency 200-hPa height Z and envelope storm track E (see text for details). No. PCs gives the number of PCs used, Tot sq cov gives the total squared covariance in m2, Sq cov is the squared covariance for each mode in m2, Cov exp is the percentage of squared covariance explained by each mode, Z var exp is the percentage of height variance explained, Z cor is the correlation (×100) between the PC of the Z and MCA mode, and E cor is the correlation (×100) between the PC of the E and MCA modes.", null, "One robust result from Table 1 is that the first 3 MCA modes dominate at all truncations, explaining at least 75% of the total squared covariance. At least five PCs are required, however, to represent these three modes in a consistent fashion.\n\nFigures 1 –3 compare the MCA patterns of Z and E for truncations of 5 and 10 PCs for the leading three modes. All the patterns shown are heterogeneous correlation functions (Bretherton et al. 1992), so that the maps of E give the covariances of this envelope field with the (standardized) Z time series associated with the given mode, while the height patterns give the covariance of Z with the (standardized) E time series associated with the given mode.\n\nFig. 1.\n\nHeterogeneous covariance maps for first MCA mode with truncation of 10 PCs for (a) the low-frequency height field and (b) the envelope rms of high-frequency meridional wind. (c),(d) The same results for 5 PC truncation. The covariance of each field is shown with respect to the unit of standard deviation of the other fields’ MCA temporal coefficient. The contour interval is 10 m in (a),(c) and 0.2 m s−1 in (b),(d).\n\nFig. 1.\n\nHeterogeneous covariance maps for first MCA mode with truncation of 10 PCs for (a) the low-frequency height field and (b) the envelope rms of high-frequency meridional wind. (c),(d) The same results for 5 PC truncation. The covariance of each field is shown with respect to the unit of standard deviation of the other fields’ MCA temporal coefficient. The contour interval is 10 m in (a),(c) and 0.2 m s−1 in (b),(d).\n\nFig. 3.\n\nHeterogeneous covariance maps for the third MCA mode with truncation of 10 PCs for (a) the low-frequency height field and (b) the envelope rms of high-frequency meridional wind. (c),(d) The same results for 5 PC truncation. The covariance of each field is shown with respect to the unit of standard deviation of the other fields’ MCA temporal coefficient. The contour interval is 10 m in (a),(c) and 0.2 m s−1 in (b),(d).\n\nFig. 3.\n\nHeterogeneous covariance maps for the third MCA mode with truncation of 10 PCs for (a) the low-frequency height field and (b) the envelope rms of high-frequency meridional wind. (c),(d) The same results for 5 PC truncation. The covariance of each field is shown with respect to the unit of standard deviation of the other fields’ MCA temporal coefficient. The contour interval is 10 m in (a),(c) and 0.2 m s−1 in (b),(d).\n\nThe leading MCA pattern shows a strong high in the North Pacific and a trough over northwestern North America (Figs. 1a,c). The associated northward shift in the storm track (seen in the envelope of υrms in Figs. 1b,d) is consistent with the shift in the jet implied by Figs. 1a,c.\n\nThere is agreement in the overall pattern between the two representations based on a different truncation. This is also true for the second mode (Fig. 2), which shows a North Atlantic mode combined with a low over western North America. Here we see that the weakened westerlies (anomalous easterlies) in the North Atlantic are consistent with the reduced storm track (envelope) activity, while farther south the situation is reversed. In the Pacific, the low over western North America (high in the North Pacific) is also associated with consistent storm track changes.\n\nFig. 2.\n\nHeterogeneous covariance maps for the second MCA mode with truncation of 10 PCs for (a) the low-frequency height field and (b) the envelope rms of high-frequency meridional wind. (c),(d) The same results for 5 PC truncation. The covariance of each field is shown with respect to the unit of standard deviation of the other fields’ MCA temporal coefficient. The contour interval is 10 m in (a),(c) and 0.2 m s−1 in (b),(d).\n\nFig. 2.\n\nHeterogeneous covariance maps for the second MCA mode with truncation of 10 PCs for (a) the low-frequency height field and (b) the envelope rms of high-frequency meridional wind. (c),(d) The same results for 5 PC truncation. The covariance of each field is shown with respect to the unit of standard deviation of the other fields’ MCA temporal coefficient. The contour interval is 10 m in (a),(c) and 0.2 m s−1 in (b),(d).\n\nThe third MCA pattern shows a sort of asymmetric annular mode, with enhanced height gradients in the Pacific (and a low over Alaska), consistent with enhanced storm track activity in this region. This configuration holds also in the North Atlantic, although the height gradient and envelope function anomalies are fairly weak in the five PC truncation.\n\n## 3. Cluster analysis and significance\n\n### a. Truncation and algorithm\n\nWe apply the cluster analysis directly to the set of Z time series associated with the first M MCA modes, where the dimension M is to be chosen. (Each atmospheric state is represented by a point in a M-dimensional space.) The cluster analysis consists of a partitioning algorithm that, given a prespecified number k of clusters, iteratively assigns all points to one of k groups (clusters) so that the ratio of the within-cluster variance to the total variance is minimized (Michelangeli et al. 1995). No quasi-stationary filtering was applied (see SCM). Once each state is identified as belonging to a particular cluster, corresponding composite maps of Z and E are obtained using the full original gridpoint representation of the filtered datasets. The composites are presented as anomalies from the full climatology for all fields.\n\nTo apply the cluster analysis we must first decide the dimension M and the number N of PCs (i.e., the truncation) to be used to define the MCA modes. The results of the previous section suggest that only three distinct MCA modes explain the covariance of the storm track envelope field with the low-frequency height field. A minimum truncation of 5 PCs for each field is necessary to capture these MCA modes, suggesting that the truncation N be at least 5. We find consistent cluster analysis results for a range of truncations N ≥ 5 and dimension M ≥ 3.\n\nFigures 4 –7 compare the composite fields of Z and E for 4 clusters (k = 4) for the choices of (N = 5 and M = 3), (N = 5 and M = 5), (N = 10 and M = 5), and (N = 10 and M = 10). The height composites are indicated by the contours, the envelope υrms by the shading. The Pacific trough, Alaska ridge, and Arctic high patterns of SCM show reasonably good agreement across the various choices of truncation and dimension, with the envelope field consistent with the height gradient anomalies.\n\nFig. 4.\n\nPacific trough cluster composites using (a) 10 MCA coefficients computed from 10 PCs, (b) 5 MCA coefficients computed from 10 PCs, (c) 5 MCA coefficients computed from 5 PCs, and (d) 3 MCA coefficients computed from 5 PCs. The contours show the 200-hPa height field (contour interval of 30 m), shading the 200-hPa envelope rms of high-frequency meridional wind (interval of 0.3 m s−1).\n\nFig. 4.\n\nPacific trough cluster composites using (a) 10 MCA coefficients computed from 10 PCs, (b) 5 MCA coefficients computed from 10 PCs, (c) 5 MCA coefficients computed from 5 PCs, and (d) 3 MCA coefficients computed from 5 PCs. The contours show the 200-hPa height field (contour interval of 30 m), shading the 200-hPa envelope rms of high-frequency meridional wind (interval of 0.3 m s−1).\n\nFig. 7.\n\nArctic low cluster composites using (a) 10 MCA coefficients computed from 10 PCs, (b) 5 MCA coefficients computed from 10 PCs, (c) 5 MCA coefficients computed from 5 PCs, and (d) 3 MCA coefficients computed from 5 PCs. The contours show the 200-hPa height field (contour interval of 30 m), shading the 200-hPa envelope rms of high-frequency meridional wind (interval of 0.5 m s−1).\n\nFig. 7.\n\nArctic low cluster composites using (a) 10 MCA coefficients computed from 10 PCs, (b) 5 MCA coefficients computed from 10 PCs, (c) 5 MCA coefficients computed from 5 PCs, and (d) 3 MCA coefficients computed from 5 PCs. The contours show the 200-hPa height field (contour interval of 30 m), shading the 200-hPa envelope rms of high-frequency meridional wind (interval of 0.5 m s−1).\n\nThe Arctic low pattern shows somewhat less consistency for the (N = 10 and M = 5) case. This cluster was also found to be less robust in SCM. Repeating the cluster analysis with only k = 3 clusters, we find good agreement between all choices of N and M, and with the results of SCM.\n\n### b. Significance\n\nCluster analysis seeks to identify regions in the reduced space whose probability density is higher than would be expected from a suitable background (a multinormal distribution). Although PCs are linearly uncorrelated (orthogonal in time), a “significant” cluster implies that they are not statistically independent. Significance testing is carried out by generating a large number of independent synthetic time series for each PC. Using the approach we describe below, we are assured that each synthetic time series approximately captures the temporal correlation structure of the corresponding real PC, yet they are statistically independent.\n\nThis approach, designed for clusters computed from PCs, would not be appropriate if applied directly to an MCA-based cluster calculation, since the MCA time series for different modes are not linearly uncorrelated. However, if synthetic PCs are generated as described below, and the MCA and cluster analysis repeated using these PCs, the results are appropriate for assessing significance.\n\nChristiansen (2007) pointed out that the temporal correlation structure of any time series can be captured in a synthetic time series by computing the Fourier harmonic coefficients for all frequencies, retaining the original amplitudes, but randomizing the phases for each frequency. This random-phase approach preserves the spectrum. Since (for an infinite time series) the spectrum is the Fourier transform of the temporal autocovariance function, the synthetic time series should retain the original temporal correlation information.\n\nThis approach needs modification when applied to winter daily data, consisting of discontinuous segments. Write the principal components for Z, Py,d (where the index y gives the year and the index d gives the day), as\n\nwhere 1 ≤ y ≤ 54, 1 ≤ d ≤ 92, the seasonal mean S depends only on year y, and the series F has zero mean over all d for any y. The random-phase approach is implemented on the series S and F separately. Here Fy,d represents a sample (of size 54) of daily time series of zero mean; the spectrum is computed for each sample and then averaged over all available samples. A random-phase surrogate d is computed by summing the harmonic coefficients computed using the specified amplitudes with random phases. The series S, however, consists of only one sample of 54 seasonal means. To avoid sampling errors, its spectrum is approximated by a white-noise spectrum, whose constant value is just the average over all frequencies. A surrogate Ŝy is constructed again using random phases.\n\nA synthetic time series corresponding to Py,d is formed using a randomly chosen surrogate Ŝy to which we add 54 randomly chosen surrogates d. This is repeated for each of the PCs to form a synthetic dataset of the same size as the original observed dataset. The cluster analysis is repeated for each of 100 such datasets; the number of datasets for which the variance ratio is less than the observed (synthetic data more clustered than observed data), once divided by 100, gives the significance level.\n\nSignificance results are shown in Table 2 for a range of k and a range of M, the number of MCA modes used in the clustering. In each case the truncation N = M. For 3 or 4 clusters, the corresponding p value is 0.03 or below3 for all M. Note that this significance test does not measure the degree to which the underlying distribution departs from a multinormal one (Christiansen 2007).\n\nTable 2.\n\nStatistical significance of clusters using MCA coefficients as a function of the dimension of the space M (number of MCA modes used), and the number of clusters k. The numbers reported are the p values, computed as the percentage of synthetic time series with greater clustering (lower variance ratio) than the reanalysis time series. In all cases, the number of PCs N = M.", null, "## 4. Conclusions and discussion\n\n• For wintertime in the Pacific–North America region, three MCA modes capture over 75% of the squared covariance between the low-frequency 200-hPa height field and the envelope fields describing the low-frequency variation of the rms of high-frequency meridional wind. This result is robust to changes in the PC truncation of the fields in the analysis as long as at least 5 PCs are retained.\n\n• Cluster analysis carried out for a range of PC truncations and number of MCA modes identifies essentially the same preferred patterns identified in SCM (for the same observational dataset) on the basis of the low-frequency flow. These states are thus characterized by a strong eddy–mean flow relationship.\n\n• An analysis of significance improved over that of SCM shows that 3 or 4 clusters are significant with a p value of 0.03 or less, independent of the number of modes kept in the cluster analysis. This is based on the full input dataset. The degree to which the underlying distribution departs from a Gaussian in not measured in this analysis.\n\nFor the region and season presented here, inclusion of the transient eddy configuration did not substantially modify the preferred patterns suggested by the original cluster analysis using principal components. However in general this will not be the case.\n\nFig. 5.\n\nAlaska ridge cluster composites using (a) 10 MCA coefficients computed from 10 PCs, (b) 5 MCA coefficients computed from 10 PCs, (c) 5 MCA coefficients computed from 5 PCs, and (d) 3 MCA coefficients computed from 5 PCs. The contours show the 200-hPa height field (contour interval of 30 m), shading the 200-hPa envelope rms of high-frequency meridional wind (interval of 0.3 m s−1).\n\nFig. 5.\n\nAlaska ridge cluster composites using (a) 10 MCA coefficients computed from 10 PCs, (b) 5 MCA coefficients computed from 10 PCs, (c) 5 MCA coefficients computed from 5 PCs, and (d) 3 MCA coefficients computed from 5 PCs. The contours show the 200-hPa height field (contour interval of 30 m), shading the 200-hPa envelope rms of high-frequency meridional wind (interval of 0.3 m s−1).\n\nFig. 6.\n\nArctic high cluster composites using (a) 10 MCA coefficients computed from 10 PCs, (b) 5 MCA coefficients computed from 10 PCs, (c) 5 MCA coefficients computed from 5 PCs, and (d) 3 MCA coefficients computed from 5 PCs. The contours show the 200-hPa height field (contour interval of 30 m), shading the 200-hPa envelope rms of high-frequency meridional wind (interval of 0.5 m s−1).\n\nFig. 6.\n\nArctic high cluster composites using (a) 10 MCA coefficients computed from 10 PCs, (b) 5 MCA coefficients computed from 10 PCs, (c) 5 MCA coefficients computed from 5 PCs, and (d) 3 MCA coefficients computed from 5 PCs. The contours show the 200-hPa height field (contour interval of 30 m), shading the 200-hPa envelope rms of high-frequency meridional wind (interval of 0.5 m s−1).\n\n## Acknowledgments\n\nThe author would like to acknowledge extensive discussions with Dr. Tim DelSole. This work was supported by the National Science Foundation under Grant ATM-03-32910, the National Aeronautics and Space Administration under Grant NNG-04-GG46G, and the National Oceanic and Atmospheric Administration under Grant NA04-OAR-4310034.\n\n## REFERENCES\n\nREFERENCES\nArnott\n,\nJ. M.\n,\nJ. L.\nEvans\n, and\nR.\nChiaromonte\n,\n2004\n:\nCharacterization of extratropical transition using cluster analysis.\nMon. Wea. Rev.\n,\n132\n,\n2916\n2937\n.\nAyrault\n,\nF.\n,\nF.\nLalaurette\n,\nA.\nJoly\n, and\nC.\nLoo\n,\n1995\n:\nNorth Atlantic ultra high variability.\nTellus\n,\n47A\n,\n671\n696\n.\nBlackmon\n,\nM. L.\n,\n1976\n:\nA climatological spectral study of the 500 mb geopotential height of the Northern Hemisphere.\nJ. Atmos. Sci.\n,\n33\n,\n1607\n1623\n.\nBretherton\n,\nC. S.\n,\nC.\nSmith\n, and\nJ. M.\nWallace\n,\n1992\n:\nComparison of methods for finding coupled patterns in climate data.\nJ. Climate\n,\n5\n,\n541\n560\n.\nCassou\n,\nC.\n,\nL.\nTerray\n, and\nA.\nPhillips\n,\n2005\n:\nTropical Atlantic influence on European heat waves.\nJ. Climate\n,\n18\n,\n2805\n2811\n.\nChristiansen\n,\nB.\n,\n2005\n:\nThe shortcomings of nonlinear principal component analysis in identifying circulation regimes.\nJ. Climate\n,\n18\n,\n4815\n4832\n.\nChristiansen\n,\nB.\n,\n2007\n:\nAtmospheric circulation regimes: Can cluster analysis provide the number?\nJ. Climate\n,\n20\n,\n2229\n2250\n.\nMichelangeli\n,\nP-A.\n,\nR.\nVautard\n, and\nB.\nLegras\n,\n1995\n:\nWeather regimes: Recurrence and quasi stationarity.\nJ. Atmos. Sci.\n,\n52\n,\n1237\n1256\n.\nMonahan\n,\nA. H.\n,\nL.\nPandolfo\n, and\nG.\nFlato\n,\n2001\n:\nThe preferred structure of variability of the northern hemisphere atmospheric circulation.\nGeophys. Res. Lett.\n,\n28\n,\n1019\n1028\n.\nNakamura\n,\nH.\n, and\nJ. M.\nWallace\n,\n1990\n:\nObserved changes in baroclinic wave activity during the life cycles of low-frequency circulation anomalies.\nJ. Atmos. Sci.\n,\n47\n,\n1100\n1116\n.\nNakamura\n,\nH.\n,\nM.\nNakamura\n, and\nJ. L.\nAnderson\n,\n1997\n:\nThe role of high- and low-frequency dynamics in blocking formation.\nMon. Wea. Rev.\n,\n125\n,\n2074\n2093\n.\nReinhold\n,\nB. B.\n, and\nR. T.\nPierrehumbert\n,\n1982\n:\nDynamics of weather regimes: Quasi-stationary waves and blocking.\nMon. Wea. Rev.\n,\n110\n,\n1105\n1145\n.\nSimmons\n,\nA. J.\n, and\nB. J.\nHoskins\n,\n1978\n:\nThe life cycles of some nonlinear baroclinic waves.\nJ. Atmos. Sci.\n,\n35\n,\n414\n432\n.\nSmyth\n,\nP.\n,\nK.\nIde\n, and\nM.\nGhil\n,\n1999\n:\nMultiple regimes in Northern-Hemisphere height fields via mixture model clustering.\nJ. Atmos. Sci.\n,\n56\n,\n3704\n3723\n.\nStan\n,\nC.\n, and\nD. M.\nStraus\n,\n2007\n:\nIs blocking a circulation regime?\nMon. Wea. Rev.\n,\n135\n,\n2406\n2413\n.\nStraus\n,\nD. M.\n,\nS.\nCorti\n, and\nF.\nMolteni\n,\n2007\n:\nCirculation regimes: Chaotic variability versus SST-forced predictability.\nJ. Climate\n,\n20\n,\n2251\n2272\n.\nVautard\n,\nR.\n, and\nB.\nLegras\n,\n1988\n:\nOn the source of midlatitude low-frequency variability. Part II: Nonlinear equilibration of weather regimes.\nJ. Atmos. Sci.\n,\n45\n,\n2845\n2867\n.\nWallace\n,\nJ. M.\n,\nC.\nSmith\n, and\nC. S.\nBretherton\n,\n1992\n:\nSingular value decomposition of wintertime sea surface temperature and 500-mb height anomalies.\nJ. Climate\n,\n5\n,\n561\n576\n.\n\n## Footnotes\n\nCorresponding author address: David M. Straus, Center for Ocean–Land–Atmosphere Studies, 4041 Powder Mill Rd., Suite 302, Calverton, MD 20705. Email: [email protected]\n\n1\n\nMCA is often called singular value decomposition, which describes the mathematical technique used to maximize covariance between two sets of fields.\n\n2\n\nIf the Z field were to dominate the MCA, which may occur if the envelope field were to have a completely white spectrum, the matrix that describes the rotation of the Z PCs to give the new MCA coordinates would look very much like the identity matrix, at least for the lowest modes (upper-left-hand corner). We have confirmed that this is not the case here, so that the rotation matrix looks very distinct from the identity matrix.\n\n3\n\nThe high significance for k = 6 clusters is also found in SCM. In general the significance increases with k, but this is offset by a decrease in reproducibility, as assessed by comparing results from half-length samples. This issue is discussed in detail in SCM." ]
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https://trac.ffmpeg.org/ticket/1108
[ "Opened 10 years ago\n\nClosed 10 years ago\n\n# Always using \"default\" memory handling from mem.c against own implementation of memory allocator and deallocator.\n\nReported by: Owned by: mato Michael Niedermayer normal avutil git-master memory, AVDictionary, alloc, free no no\n\n## Description\n\nIn library avutil function av_free is used and other functions from mem.c. If I have my own memory handling functions, in libavutil always default is used instated of mine. I had situation when for allocation my function was used, however for freeing - default av_free was used from file dict.c from libavutil library. It can lead to crashes.\n\n### comment:1 by Carl Eugen Hoyos, 10 years ago\n\nIs this a regression?\nIs this also reproducible with current git head?\n\nCould you add source code that crashes or at least a backtrace of the crash?\n\n### comment:2 by Michael Niedermayer, 10 years ago\n\nPriority: important → normal 0.10 → git-master\n\nWhy do you want to use your own memory allocation functions with AVDict ?\nIam asking so we can understand if there is a use case for this or if we just need to document better that this isnt supported\n\n### follow-up:  5 comment:3 by mato, 10 years ago\n\nHi,\n\n• Is this a regression? - i think it is not regression, problem exist in ffmpeg 0.6.2 but in other file: avstring.c\n• Is this also reproducible with current git head? - my project use 0.10 ffmpeg version and if i try compile with head not all is compatible\n• Could you add source code that crashes or at least a backtrace of the crash?\n\nYes, of course:\nThere is gdb backtrace:\ngdb Backtrace:\n\n0xb7fe1424 in kernel_vsyscall ()\n(gdb) bt\n#0 0xb7fe1424 in\nkernel_vsyscall ()\n#1 0xb75dce71 in raise (sig=6) at ../nptl/sysdeps/unix/sysv/linux/raise.c:64\n#2 0xb75e034e in abort () at abort.c:92\n#3 0xb76147f7 in libc_message (do_abort=2, fmt=0xb76edc0c \"* glibc detected * %s: %s: 0x%s *\\n\") at ../sysdeps/unix/sysv/linux/libc_fatal.c:189\n#4 0xb761ebe1 in malloc_printerr (action=<value optimized out>, str=<value optimized out>, ptr=0x807d900) at malloc.c:6283\n#5 0xb762050b in _int_free (av=<value optimized out>, p=0x807d8f8) at malloc.c:4795\n#6 0xb762369d in\nlibc_free (mem=0x807d900) at malloc.c:3738\n#7 0xb7749584 in av_dict_free () from /media/truecrypt1/work/trunk2_3/src/dlna/target/linux/bin/dms_smm/ffmpeg_libs/libavutil.so.51\n#8 0x0807d760 in ?? ()\nBacktrace stopped: previous frame inner to this frame (corrupt stack?)\n\nAnd I think more intersting valgrind log for version 0.10:\n\nvalgrind log:\n==18131== Invalid free() / delete / delete[]\n==18131== at 0x4025BF0: free (vg_replace_malloc.c:366)\n==18131== by 0x48B6583: av_dict_free (in /media/truecrypt1/work/trunk2_3/src/dlna/target/linux/bin/dms_smm/ffmpeg_libs/libavutil.so.51) You can se that free mettod is system one\n==18131== Address 0x4ccad00 is 16 bytes inside a block of size 41 alloc'd\n==18131== at 0x402517B: memalign (vg_replace_malloc.c:581)\n==18131== by 0x8054C00: alloc_ram (av_mem.c:382)\n==18131== by 0x805545C: av_malloc (av_mem.c:699)\n==18131== by 0x8055554: av_realloc (av_mem.c:740)\n==18131== by 0x8054669: av_realloc_f (av_mem.c:186)\n==18131== by 0x40AEAFC: dyn_buf_write (in /media/truecrypt1/work/trunk2_3/src/dlna/target/linux/bin/dms_smm/ffmpeg_libs/libavformat.so.53)\nbut here function from libavformat realloc use memory handling from av_mem.c (own memory implementation)\n==18131==\n==18131== Invalid free() / delete / delete[]\n==18131== at 0x4025BF0: free (vg_replace_malloc.c:366)\n==18131== by 0x48B6592: av_dict_free (in /media/truecrypt1/work/trunk2_3/src/dlna/target/linux/bin/dms_smm/ffmpeg_libs/libavutil.so.51)\n==18131== Address 0x4ccb3f0 is 16 bytes inside a block of size 48 alloc'd\n==18131== at 0x402517B: memalign (vg_replace_malloc.c:581)\n==18131== by 0x8054C00: alloc_ram (av_mem.c:382)\n==18131== by 0x805545C: av_malloc (av_mem.c:699)\n==18131== by 0x8055554: av_realloc (av_mem.c:740)\n==18131== by 0x8054669: av_realloc_f (av_mem.c:186)\n==18131== by 0x40AEAFC: dyn_buf_write (in /media/truecrypt1/work/trunk2_3/src/dlna/target/linux/bin/dms_smm/ffmpeg_libs/libavformat.so.53)\n==18131==\n\nand valgrind log from 0.6.2 version:\n==653== Invalid read of size 4\n==653== at 0x4027D42: memcpy (mc_replace_strmem.c:635)\n==653== by 0x8052CAE: dlna_memcpy (peer_memory.c:54)\n==653== by 0x8055585: av_free (av_mem.c:779)\n==653== by 0x40B9627: av_metadata_free (metadata.c:103) and here av_metatdata_free use av_mem.c (ovn implementation)\n==653== by 0x80501E7: SMDM_server_handle_client (mpe_server.c:338)\n==653== by 0x804FE3D: main (mpe_server.c:208)\n==653== Address 0x4c3facc is 4 bytes before a block of size 16 alloc'd\n==653== at 0x402517B: memalign (vg_replace_malloc.c:581)\n==653== by 0x40251D8: posix_memalign (vg_replace_malloc.c:709)\n==653== by 0x47527DA: av_malloc (mem.c:83)\n==653== by 0x4750743: av_d2str (avstring.c:96)\nhere av_d2str from avstring.c alloc memory from system\n==653== by 0x40ED783: av_open_input_stream (utils.c:458)\n==653== by 0x40F1ED6: av_open_input_file (utils.c:612)\n==653== by 0x804D27F: SMDM_parser_mpe_query (mpe_parse.c:312)\n==653== by 0x80501E7: SMDM_server_handle_client (mpe_server.c:338)\n==653==\n==653== Invalid free() / delete / delete[]\n==653== at 0x4025BF0: free (vg_replace_malloc.c:366)\n==653== by 0x8054CF8: free_ram (av_mem.c:477)\n==653== by 0x80555B9: av_free (av_mem.c:787)\n==653== by 0x80501E7: SMDM_server_handle_client (mpe_server.c:338)\n==653== by 0x804FE3D: main (mpe_server.c:208)\n==653== Address 0x4c3fac0 is 16 bytes before a block of size 16 alloc'd\n==653== at 0x402517B: memalign (vg_replace_malloc.c:581)\n==653== by 0x40251D8: posix_memalign (vg_replace_malloc.c:709)\n==653== by 0x47527DA: av_malloc (mem.c:83)\n==653== by 0x4750743: av_d2str (avstring.c:96)\n==653== by 0x40ED783: av_open_input_stream (utils.c:458)\n==653== by 0x40F1ED6: av_open_input_file (utils.c:612)\n==653== by 0x804D27F: SMDM_parser_mpe_query (mpe_parse.c:312)\n==653== by 0x80501E7: SMDM_server_handle_client (mpe_server.c:338)\n\nAccording to description i mem.c file: line 64\n\n/* You can redefine av_malloc and av_free in your project to use your\n\nmemory allocator. You do not need to suppress this file because the\nlinker will do it automatically. */\n\nI redefined memory handling functions. But You can see in logs that sometimes allocation is called from other place than free. Unfortunatly all av_malloc and av_free functions that makes a problem are called from inside ffmpeg.\n\nBest Regards\nMarcin Tomczyk\n\n### comment:4 by mato, 10 years ago\n\nIf You analize log below You can see that dyn_buf_write alloc memory in libavformat and\nav_dict_free free in libavutil. If I redefine av_malloc and av_free dyn_buf_write use my implementation but dyn_buf_write use default implementation. That is a problem.\n\n==18131== Invalid free() / delete / delete[]\n==18131== at 0x4025BF0: free (vg_replace_malloc.c:366)\n==18131== by 0x48B6583: av_dict_free (in /media/truecrypt1/work/trunk2_3/src/dlna/target/linux/bin/dms_smm/ffmpeg_libs/libavutil.so.51) You can se that free mettod is system one\n==18131== Address 0x4ccad00 is 16 bytes inside a block of size 41 alloc'd\n==18131== at 0x402517B: memalign (vg_replace_malloc.c:581)\n==18131== by 0x8054C00: alloc_ram (av_mem.c:382)\n==18131== by 0x805545C: av_malloc (av_mem.c:699)\n==18131== by 0x8055554: av_realloc (av_mem.c:740)\n==18131== by 0x8054669: av_realloc_f (av_mem.c:186)\n==18131== by 0x40AEAFC: dyn_buf_write (in /media/truecrypt1/work/trunk2_3/src/dlna/target/linux/bin/dms_smm/ffmpeg_libs/libavformat.so.53) but here function from libavformat realloc use memory handling from av_mem.c (own memory implementation)\n==18131==\n==18131== Invalid free() / delete / delete[]\n==18131== at 0x4025BF0: free (vg_replace_malloc.c:366)\n==18131== by 0x48B6592: av_dict_free (in /media/truecrypt1/work/trunk2_3/src/dlna/target/linux/bin/dms_smm/ffmpeg_libs/libavutil.so.51)\n==18131== Address 0x4ccb3f0 is 16 bytes inside a block of size 48 alloc'd\n==18131== at 0x402517B: memalign (vg_replace_malloc.c:581)\n==18131== by 0x8054C00: alloc_ram (av_mem.c:382)\n==18131== by 0x805545C: av_malloc (av_mem.c:699)\n==18131== by 0x8055554: av_realloc (av_mem.c:740)\n==18131== by 0x8054669: av_realloc_f (av_mem.c:186)\n==18131== by 0x40AEAFC: dyn_buf_write (in /media/truecrypt1/work/trunk2_3/src/dlna/target/linux/bin/dms_smm/ffmpeg_libs/libavformat.so.53)\n\nnext log:" ]
[ null ]
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https://www.physicsforums.com/threads/rotating-ideal-gas.232308/
[ "# Rotating Ideal Gas\n\n## Homework Statement\n\nAn ideal gas of N molecules of mass m is contained in a cylinder of length\nL and radius R. The cylindrical container is rotating about its axis at an\nangular velocity $$\\omega$$, and is at equilibrium with temperature T.\n\nWrite down the energy for single particle states, and the Boltzmann distribution for the gas in the rotating cylinder given it experiences a centrifugal potential energy V(r)\n\n## Homework Equations\n\nthe energy is just E = $$\\frac{p^2}{2m} + V(r)$$\n\nI'm not entirely sure about the Boltzmann distribution. In my notes it states that $$Z_1 = \\int \\frac{d^3 rd^3p}{h^3} e^{-\\beta E}$$\n\nHowever here I get confused. I was pretty sure that Z is the partition function of the gas (which it has been so far), but a few lines further in the notes it calls this the Boltzmann distribution. As far as I was aware these aren't interchangeable names, and are completely separate things within the frame work of statistics.\n\nThe question goes onto to ask me to split the Boltzmann distribution into a translational and a interation part, which I know can be done with the above equation (another reason I'm hedging my bets on using this equation for the Boltzmann distribution). Whilst this is fine in theory I'm a little concerned about doing this in practice. I know one of the parts is an integral over 3d space with the kinetic part of the energy and the other bit is the integral over what seems to be 3d momentum with the potential part (I think I may have got the combinations back to front, but I'll check that - thats not the bit I'm confused with). My problem with this is the 3d momentum integral. I have never seen how to do this (or even seen anything like it before!) so could you suggest how I'd go about solving this bit. Is it anything like the the 3d integral over space in which you have various factors to add to it, or does it work completely differently?\n\nLast edited:\n\n$$e^{-\\beta E}$$\n$$\\int_{-\\infty}^{\\infty}dx e^{-\\alpha x^2}=\\sqrt{\\frac{\\pi}{\\alpha}}$$" ]
[ null ]
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https://nickadamsinamerica.com/missing-angles-in-triangles-worksheet-answers-pdf/e9c25ed5ed351a5a5e9e48e4b5135894-2/
[ "# E9c25ed5ed351a5a5e9e48e4b5135894\n\nBy . Worksheet. At Saturday, October 09th 2021, 23:19:04 PM.\n\nThese multiplication worksheets may be configured for either single or multiple digit horizontal problems. The factors may be selected to be positive, negative or mixed numbers for these multiplication worksheets. You may vary the numbers of multiplication problems on the multiplication worksheets from 12 to 30. These multiplication worksheets are appropriate for Kindergarten, 1st Grade, 2nd Grade, 3rd Grade, 4th Grade, and 5th Grade.\n\nThese addition worksheets will produce 12 vertical or horizontal addition problems using dot figures to represent the numbers. You may select the numbers for the addition worksheets to be used from 0 to 10.\n\nThese Radical Worksheets will produce problems for multiplying radical expressions. You may select the difficulty for each expression. These Radical Worksheets are a good resource for students in the 5th Grade through the 8th Grade.", null, "", null, "Angle Relationships Maze Finding Angle Measures Finding Angle Measures Angle Relationships Angle Relationships Worksheet", null, "Similar Figures Middle School Math Geometry Math Work Geometry High School", null, "Missing Angles In Triangles Triangles Angle Sum Theorem Angles Math Teaching Geometry Triangle Math", null, "Triangle Congruence Proofs Drag And Drop Activity Activities Geometry Activities Resource Classroom\n\n### Gallery of Missing Angles In Triangles Worksheet Answers Pdf\n\n1 star 2 stars 3 stars 4 stars 5 stars\n\nAny content, trademark/s, or other material that might be found on this site that is not this site property remains the copyright of its respective owner/s." ]
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http://www.cad.zju.edu.cn/home/zhx/csmath/doku.php?id=2011:pde
[ "# CSMATH\n\n2011:pde\n\n## Introduction to Partial Differential Equation\n\nPartial Differential Equation is an important mathematical tool for computer science. This one month short course is designed to give first year Ph.D. students a thorough grounding in the methodologies, technologies, mathematics and algorithms currently needed by people who are doing research related to PDE. The topics of the course draw mainly from types of image and video processing, from geometrical modeling, and from statistical algorithmics.\n\nStudents entering the class should have a pre-existing working knowledge of fundamental mathematics and algorithms, though the class has been designed to allow students with a strong numerate background to catch up and fully participate.\n\n## Schedule\n\nTopic Date Slides note\nIntroduction 2011.04.28 introduction =>\nPoisson Equation 2011.05.05 Poisson Equation =>\nCourse talk\n(Xuezhen Huang)\n2011.05.05 Diffusion Curves: A Vector Representation for Smooth-Shaded Images\nLevel set (I) 2011.05.12 slides =>\nLevel set (II) 2011.05.19 slides =>\n\n## Text books\n\n1. Arieh Lserles. A first course in the numerical analysis of differential equations. (微分方程数值分析基础教程, 清华大学出版社)\n2. 余德浩, 汤华中. 微分方程数值解法. 科学出版社.\n3. Stanley Osher, Ron Fedkiw. Level Set Methods and Dynamic Implicit Surfaces.", null, "" ]
[ null, "http://www.cad.zju.edu.cn/home/zhx/csmath/lib/exe/indexer.php", null ]
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https://codereview.stackexchange.com/questions/134380/pairwise-programming-challenge-free-code-camp
[ "# Pairwise programming challenge - Free code camp\n\nChallenge:\n\nGiven an array arr, find element pairs whose sum equal the second argument arg and return the sum of their indices.\n\nIf multiple pairs, that have the same numeric elements but different indices are possible, return the smallest sum of indices. Once an element has been used, it cannot be reused to pair with another.\n\nFor example pairwise([7, 9, 11, 13, 15], 20) returns 6.\n\nMy code (works, passes all validations) is the below:\n\nfunction pairwise(list, sum) {\n// noprotect\nvar innerCounter,\nouterCounter,\nused = {},\nreturnValue = 0;\nfor( outerCounter = 0 ; outerCounter < list.length-1 ; outerCounter++){\nif( used[outerCounter] )\ncontinue;\nfor( innerCounter = outerCounter+1 ; innerCounter < list.length ; innerCounter++){\nif( used[innerCounter] )\ncontinue;\nif( list[innerCounter] + list[outerCounter] == sum ){\nused[innerCounter] = true;\nreturnValue = returnValue + innerCounter + outerCounter;\nbreak;\n}\n}\n}\n\nreturn returnValue;\n}\n\n\nI had some internal debate on sum as a parameter name, since we also return a sum. But I stuck to returnValue for the return value and sum for the parameter. Any feedback on this is welcome, from style to algorithm.\n\n• Are the elements of list necessarily distinct? – 200_success Jul 9 '16 at 21:26\n• @200_success Not at all. In fact [0,0,0,1,1],1 is a test case. – konijn Jul 10 '16 at 23:17\n\n## 2 Answers\n\nMy only comments would be stylistic so may be more a matter of opinion. Am open to any counter comments.\n\nWhitespace\n\nAll your code is bunched together and is not as easy to read as it could be.\n\nSingle Statement if braces\n\nIt is often beneficial and consistent to always have braces on if and other blocks even if there is only a single statement. It's easy to introduce undesirable functionality if you come back to the code later and comment/add lines e.g.\n\nif(someValue)\n//doThis();\nalwaysDoThis();\n\n\nShorthand arithmetic\n\nreturnValue = returnValue + innerCounter + outerCounter;\n\n\ncan be shortened to\n\nreturnValue += innerCounter + outerCounter;\n\n1. Algorithm\n\nYour code proposes a nested loop, which may indicate a high time complexity. I believe more efficient implementations exist, possibly at the price of an increased code complexity.\n\nRather than search ahead for a complement of the current value, I think it would be more interesting to store the complement of the current value, then move on - each time checking whether the current value matches an already stored complement. The first script below illustrates this approach.\n\nAlso, in such problems, it is always instructive to examine how sorting the initial set can further improve the algorithm. In this occurrence, provided the sorted values retain their initial index, the benefits are interesting: you can iterate from the start and end of the set at the same time, matching values as you go, until iterators cross over. The second script below illustrates this approach.\n\n2. Style\n\nEven though there is no block-scoping in JavaScript, you should not declare all your variables at once at the top of the function, but rather declare them when they are used. A reviewer will keep his eyes in the same code area, which will facilitate his understanding.\n\nAs pointed out by Tony, you should always use curly braces, even if it is to surround a single statement. This way, reviewers do not have to worry about comments or semi-colons insidiously invalidating the code logic.\n\nFinally, your internal debate on the naming of the returned value is quite interesting. You should indeed strive to make variables self-explanatory and readable. In your case, since you return the sum of indices, an appropriate variable name would be... sumOfIndices.\n\n3. Illustrations\n\n1. Storing complements:\n\nfunction pairwise(list, sum) {\nvar sumOfIndices = 0;\nvar complements = {};\nfor(var ii = 0; ii < list.length; ii++) {\nvar value = list[ii];\nvar complement = sum - value;\nif (typeof complements[complement] == \"undefined\") {\nvar stored = complements[value];\nif (typeof stored == \"undefined\") {\nstored = (complements[value] = []);\n}\nstored.push(ii);\n} else {\nvar matches = complements[complement];\nsumOfIndices += matches.pop() + ii;\nif (matches.length == 0) {\ndelete complements[complement];\n}\n}\n}\nreturn sumOfIndices;\n}\n\n2. Sorting values:\n\nfunction pairwise(list, sum) {\nvar sumOfIndices = 0;\nvar entries = list.map((value, index) => { return {value: value, index: index}; });\nentries.sort(function(a, b) {\nif (a.value < b.value) { return -1; }\nif (a.value > b.value) { return +1; }\nreturn 0;\n});\nvar ii = entries.length - 1;\nfor(; ii > 0; ii--) {\nif (entries[ii].value <= sum) { break; }\n}\nfor(var j = 0; j < ii; j++) {\nvar first = entries[j];\nfor(; ii > j; ii--) {\nvar second = entries[ii];\nif (first.value + second.value > sum) {\ncontinue;\n}\nif (first.value + second.value == sum) {\nsumOfIndices += first.index + second.index;\nii--; // Skip last match.\n}\nbreak;\n}\n}\nreturn sumOfIndices;\n}" ]
[ null ]
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https://fabian-voith.de/2020/10/29/subnetting-quick-easy/
[ "# Subnetting quick & easy\n\nIt’s very useful to be able to quickly perform some subnetting when given an IP address with its CIDR – it allows you to determine how many addresses there are in the corresponding network, if the IP is a broadcast IP, if another IP belongs to the same network, etc. While you might remember some “common” CIDRs like /8, /16, /24, there is also a quick and easy way how to determine the subnet for less well-known CIDRs.\n\nThere are many videos on this, but I find Sebastian Philippi’s technique the most straightforward one. As the video is only available in German and as it also omits some details which I think are helpful for a better understanding, I wanted to share this tip with you.\n\n### Composition of an IP address\n\nAs you will probably know, an IPv4 Address is composed of 4 numbers separated by a dot, with each number between 0 and 255, e.g. 192.168.178.2 or 127.0.0.1.\n\nTechnically spoken, each IP address is a 32-bit binary number.\n\nWhy 32 bits? Since each of the 4 numbers can be anything between 0 and 256, you will need up to 8 bits to represent them. 4 times an 8-bit number equals 32 bits. For example:\n\n192.168.178.2 in binary representation would be:\n\n``````192 = 1*128 + 1*64 + 0*32 + 0*16 + 0*8 + 0*4 + 0*2 + 0*1\n168 = 1*128 + 0*64 + 1*32 + 0*16 + 1*8 + 0*4 + 0*2 + 0*1\n178 = 1*128 + 0*64 + 1*32 + 1*16 + 0*8 + 0*4 + 1*2 + 0*1\n2 = 0*128 + 0*64 + 0*32 + 0*16 + 0*8 + 0*4 + 1*2 + 0*1``````\n\nIf we format this like an IP address, 192.168.178.2 becomes:\n\n11000000.10101000.10110010.00000010\n\nTip: You can easily convert this using some IP-to-Binary online converter, but in practice you will not need to do this. I just explained it here to build a better understanding of our following steps.\n\nEach IP address belongs to a network, which is indicated by a subnet mask. This bitmask is then applied with a bitwise AND to the IP address to determine the subnet of that IP. This process is called subnetting. When you look into your computer settings, you will see a subnet mask written in decimal form, just as an IP address, e.g.\n\n255.255.255.0\n\nHowever, in networking scenarios you will often see this represented as a CIDR notation instead, which stands for classless inter domain routing.\n\nNo matter if a subnet mask or a CIDR notation is used, in both cases we can see which parts of an address are “fixed” and which can be used by systems in that network. These two parts are known as the network prefix (the part which is fixed), and the host identifier (the variable part).\n\nFrom a subnet you can simply see which parts are variable. The subnet above – 255.255.255.0 – tells us that the first three of the four octets are fixed (network part) and that the last octet can be used by hosts in the network (0-256).\n\nWhen we have a CIDR, this is a bit more difficult to see, but as I said in the introduction, there is a quick and easy way how to retrieve the same information. Here is how:\n\n### Determining network prefix and host identifier\n\nRemember that each “block” of an IPv4 address is 8 bits long. If we “count” the bits with increasing numbers from left to right, we get:\n\n1-8.9-16.17-24.25-32\n\nWe can see: The first, second, third… eighth bit belong to the first octet, the 9th, 10th… 16th bit to the second, etc. With this knowledge we can see to which octet a given CIDR belongs.\n\nLet’s check with some examples:\n\nFor 192.168.178.2/7 we can see:\n7 is in 1-8, so this will be subject to subnetting in the first octet.\n\nFor 10.255.13.0/15 we can see:\n15 is in 9-16, so this will be subject to subnetting in the second octet.\n\nFor 127.15.19.1/20 we can see:\n20 is in 17-24, so this will be subject to subnetting in the third octet\n\nFor 12.13.14.15/26 we can see:\n26 is in 25-32, so this will be subject to subnetting in the fourth octet\n\n### Determining the network size\n\nLet’s continue with the third example, 127.15.19.1/20. Since we know that 20 is between 17-24, we also know that our network part and variable part are somewhere in the third octet.\n\nWe further know that the third octet starts with the 17th bit (because the bits 17-24 are in the third octet, see above), so if we count from that 17th bit all the way to the “CIDR bit”… 17, 18, 19, 20…  it means that the first 4 of the 8 bits in that third octet are fixed for the network (because the CIDR indicates the fixed network bits).\n\nWhat remains are then 8 bits minus 4 bits= 4 free bits for our host part. 2^4=16. This is the size of our network: We know that our networks will “jump” in steps of 16 in the third octet:\n\n127.15.0.0/20\n127.15.16.0/20\n127.15.32.0/20\n\nWe now know that with the /20 CIDR we have steps of 16 in the third octet. Our exemplary IP address is 127.15.19.1/20, so there is a 19 in the third octet. Above we saw that one network in a /20 CIDR will be 127.15.16.0/20 and the next 127.15.32.0/20, so the .19 lies in between these two. This means, for 127.15.19.1/20, the corresponding network address is 127.15.16.0/20.\n\nWe figured that 127.15.19.1/20 is in the network 127.15.16.0/20 and that the next network will be 127.15.32.0/20. Since the broadcast address is defined as the “last” address in a subnet, we simply can “deduct one” from that next network (127.15.32.0/20) to receive the broadcast address of our network:\n\n127.15.32.0/20 “-1” = 127.15.31.255\n\nSo, we see that 127.15.31.255 is the broadcast address for all IPs in the network 127.15.16.0/20 – the network to which our exemplary IP 127.15.19.1/20 belongs.\n\n### Determining the number of available addresses\n\nTo quickly calculate how many addresses are available in a network of that size, we can simply deduct the /20 from the 32 available bits in an IPv4 address (alternatively, from 128 bits for an IPv6 address) and take this as exponent to the basis of 2. In our example:\n\n32 bits (in an IPv4 address) – 20 (from our exemplary IP) = 12 bits\n2^12 = 4096\n\nIn our network 127.15.16.0/20 we therefore have 4096 available IP addresses.\n\n### Determining the last available host address\n\nAs a final step, we can find out the last IP address in that network that is available to be assigned to a system: It is simply one address “below” the broadcast address. This means:\n\n### Some more examples\n\nJust to make sure that you really understood what we were doing, here are some further examples:\n\nFor the IP 89.15.16.227/21 we can see that:\n\n1. /21 is in 17-24, therefore we will be working in the 3rd octet.\n2. Fixed bits in the third octet are 17, 18, 19, 20, 21, i.e. 5 fixed bits. As free bits available for our hosts we have 8-5=3 bits. We take these 3 bits as exponent for 2: 2^3=8, i.e. our networks will be separated by steps of 8 in the third octet:\n89.15.8.0/21\n89.15.16.0/21\n89.15.32.0/21\n3. Our exemplary address lies in the network 89.15.16.0/21, so for our broadcast address we “deduct 1” and get: 89.15.15.255.\n4. For the number of available addresses we get 32-21=11. Taken as exponent to the base of two we calculate: Our exemplary IP is one of 2^11=2048 available addresses in that network.\n5. The last available address in that network is “broadcast address minus 1”, therefore: 89.15.15.255 “-1” = 89.15.15.254.\n\nOk, one last example.\n\nFor the IP 192.168.170.14/28 we can see that:\n\n1. /28 is in 25-32, therefore we will be working in the 4th octet.\n2. Fixed bits in the fourth octet are 25, 26, 27, 28, i.e. 4 fixed bits. As free bits available for our hosts we have 8-4=4 bits. We take these 4 bits as exponent for 2: 2^4=16, i.e. our networks will be separated by steps of 16 in the fourth octet:\n192.168.170.0/28\n192.168.170.16/28\n192.168.170.32/28\n3. Our exemplary address lies in the network 192.168.170.0/28, so for our broadcast address we “deduct 1” and get: 192.168.169.255.\n4. For the number of available addresses we get 32-28=4. Taken as exponent to the base of two we calculate: Our exemplary IP is one of 2^4=16 available addresses in that network.\n5. The last available address in that network is “broadcast address minus 1”, therefore: 192.168.169.255 “-1” = 192.168.169.254.\n\nHappy subnetting!" ]
[ null ]
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https://www.iue.tuwien.ac.at/phd/pourfath/node69.html
[ "## 4.1.2 Boundary Conditions\n\nDIRICHLET boundary conditions are introduced at the source, drain, and gate contacts. Potentials are conveniently measured relative to the source potential. The amount of bending of the vacuum enegy level along the length of the CNT is given by", null, ", since we assume that the local electrostatic potential rigidly shifts the CNT band-structure. The conduction and valence band-edges of the CNT are given by", null, "(4.6)\n\nThe SCHOTTKY barrier heights for electrons (", null, ") and holes (", null, ") at the metal-CNT interface are given by (see Fig. 4.3)", null, "(4.7)\n\nwhere", null, "is the work function of the metal contact,", null, "is the electron affinity of the CNT, and", null, "is the band-gap of the CNT. The work function of CNT", null, "is defined as the sum of the CNT electron affinity and half of the band-gap in the bulk. Figure 4.3 shows the band diagram at the metal-CNT interface with", null, ". The work function of the CNT is assumed to be", null, ". In an intrinsic CNT (un-doped) the FERMI level of the CNT is located in the middle of the band-gap. Under these conditions, equal SCHOTTKY barrier heights for both electrons and holes are achieved. If the work functions of metal and CNT do not align, band-bending near the contact occurs and SCHOTTKY barrier heights for electrons and holes will be different. For example, if the work function of the metal contact is larger than that of the CNT (", null, ") the SCHOTTKY barrier height for holes is smaller than that for electrons. As a result, a p-type CNT-FET is achieved, where holes are the majority carriers.", null, "M. Pourfath: Numerical Study of Quantum Transport in Carbon Nanotube-Based Transistors" ]
[ null, "https://www.iue.tuwien.ac.at/phd/pourfath/img742.png", null, "https://www.iue.tuwien.ac.at/phd/pourfath/img743.png", null, "https://www.iue.tuwien.ac.at/phd/pourfath/img129.png", null, "https://www.iue.tuwien.ac.at/phd/pourfath/img130.png", null, "https://www.iue.tuwien.ac.at/phd/pourfath/img744.png", null, "https://www.iue.tuwien.ac.at/phd/pourfath/img128.png", null, "https://www.iue.tuwien.ac.at/phd/pourfath/img126.png", null, "https://www.iue.tuwien.ac.at/phd/pourfath/img138.png", null, "https://www.iue.tuwien.ac.at/phd/pourfath/img127.png", null, "https://www.iue.tuwien.ac.at/phd/pourfath/img745.png", null, "https://www.iue.tuwien.ac.at/phd/pourfath/img746.png", null, "https://www.iue.tuwien.ac.at/phd/pourfath/img747.png", null, "https://www.iue.tuwien.ac.at/phd/pourfath/img748.png", null ]
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https://adotb.xyz/theory-and-worked-examples/1-equations-and-inequalities/1-7-solving-a-system-of-linear-equations/
[ "# 1.7. Solving a system of linear equations\n\nExample question: solve\n\nSolution:\n\nOther than for solving polynomial equations , the PolySmlt2 app in our GC can also be used to solve a system of linear equations. We access the application in our calculator and go to option 2, \"Simultaneous Eqn Solver\". Set the appropriate settings (3 equations and 3 unknowns for our example) and key in our system of equations as shown in the image below. Click solve to get the solution." ]
[ null ]
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https://phys.libretexts.org/Bookshelves/College_Physics/Book%3A_College_Physics_(OpenStax)/08%3A_Linear_Momentum_and_Collisions/8.05%3A_Inelastic_Collisions_in_One_Dimension
[ "$$\\require{cancel}$$\n\n# 8.5: Inelastic Collisions in One Dimension\n\n•", null, "• OpenStax\n• General Physics at OpenStax CNX\n\nLearning Objectives\n\nBy the end of this section, you will be able to:\n\n• Define inelastic collision.\n• Explain perfectly inelastic collision.\n• Apply an understanding of collisions to sports.\n• Determine recoil velocity and loss in kinetic energy given mass and initial velocity.\n\nWe have seen that in an elastic collision, internal kinetic energy is conserved. An inelastic collision is one in which the internal kinetic energy changes (it is not conserved). This lack of conservation means that the forces between colliding objects may remove or add internal kinetic energy. Work done by internal forces may change the forms of energy within a system. For inelastic collisions, such as when colliding objects stick together, this internal work may transform some internal kinetic energy into heat transfer. Or it may convert stored energy into internal kinetic energy, such as when exploding bolts separate a satellite from its launch vehicle.\n\nDefinition: Inelastic Collisions\n\nAn inelastic collision is one in which the internal kinetic energy changes (it is not conserved).\n\nFigure $$\\PageIndex{1}$$ shows an example of an inelastic collision. Two objects that have equal masses head toward one another at equal speeds and then stick together. Their total internal kinetic energy is initially\n\n$\\dfrac{1}{2}mv^2 + \\dfrac{1}{2}mv^2 = mv^2.$\n\nThe two objects come to rest after sticking together, conserving momentum. But the internal kinetic energy is zero after the collision. A collision in which the objects stick together is sometimes called a perfectly inelastic collision because it reduces internal kinetic energy more than does any other type of inelastic collision. In fact, such a collision reduces internal kinetic energy to the minimum it can have while still conserving momentum.", null, "Figure $$\\PageIndex{1}$$: An inelastic one-dimensional two-object collision. Momentum is conserved, but internal kinetic energy is not conserved. (a) Two objects of equal mass initially head directly toward one another at the same speed. (b) The objects stick together (a perfectly inelastic collision), and so their final velocity is zero. The internal kinetic energy of the system changes in any inelastic collision and is reduced to zero in this example.\n\nDefinition: Perfectly Inelastic Collisions\n\nA collision in which the objects stick together is sometimes called “perfectly inelastic.”\n\nExample $$\\PageIndex{1}$$: Calculating Velocity and Change in Kinetic Energy - Inelastic Collision of a Puck and a Goalie\n\n1. Find the recoil velocity of a 70.0-kg ice hockey goalie, originally at rest, who catches a 0.150-kg hockey puck slapped at him at a velocity of 35.0 m/s.\n2. How much kinetic energy is lost during the collision? Assume friction between the ice and the puck-goalie system is negligible (Figure $$\\PageIndex{2}$$)", null, "Figure $$\\PageIndex{2}$$: An ice hockey goalie catches a hockey puck and recoils backward. The initial kinetic energy of the puck is almost entirely converted to thermal energy and sound in this inelastic collision.\n\nStrategy\n\nMomentum is conserved because the net external force on the puck-goalie system is zero. We can thus use conservation of momentum to find the final velocity of the puck and goalie system. Note that the initial velocity of the goalie is zero and that the final velocity of the puck and goalie are the same. Once the final velocity is found, the kinetic energies can be calculated before and after the collision and compared as requested.\n\nSolution for (a)\n\nMomentum is conserved because the net external force on the puck-goalie system is zero.\n\nConservation of momentum is\n\n$p_1 + p_2 = p'_1 + p'_2 \\nonumber$\n\nor\n\n$m_1v_1 + m_2v_2 = m_1v'_1 + m_2v'_2. \\nonumber$\n\nBecause the goalie is initially at rest, we know $$v_2 = 0.$$ Because the goalie catches the puck, the final velocities are equal, or $$v'_1 = v'_2 = v'.$$ Thus, the conservation of momentum equation simplifies to\n\n$m_1v_1 = (m_1 + m_2)v'. \\nonumber$\n\nSolving for $$v'$$ yields\n\n$v' = \\dfrac{m_1}{m_1 + m_2}v_1. \\nonumber$\n\nEntering known values in this equation, we get\n\n\\begin{align*} v' &= \\left( \\dfrac{0.150 \\, kg}{70.0 \\, kg + 0.150 \\, kg} \\right)(35.0 \\, m/s) \\\\[5pt] &= 7.48 \\times 10^{-2} m/s)^2 \\\\[5pt] &= 0.196 \\, J. \\end{align*}\n\nThe change in internal kinetic energy is thus\n\n\\begin{align*} KE'_{int} - KE_{int} &= 0.196 \\, J - 91.9 \\, J \\\\[5pt] &= -91.7 \\, J \\end{align*}\n\nDiscussion for (b)\n\nNearly all of the initial internal kinetic energy is lost in this perfectly inelastic collision. $$KE_{int}$$ is mostly converted to thermal energy and sound.\n\nDuring some collisions, the objects do not stick together and less of the internal kinetic energy is removed—such as happens in most automobile accidents. Alternatively, stored energy may be converted into internal kinetic energy during a collision. Figure $$\\PageIndex{3}$$ shows a one-dimensional example in which two carts on an air track collide, releasing potential energy from a compressed spring. Example $$\\PageIndex{2}$$ deals with data from such a collision.", null, "Figure $$\\PageIndex{3}$$: An air track is nearly frictionless, so that momentum is conserved. Motion is one-dimensional. In this collision, examined in Example $$\\PageIndex{2}$$, the potential energy of a compressed spring is released during the collision and is converted to internal kinetic energy.\n\nCollisions are particularly important in sports and the sporting and leisure industry utilizes elastic and inelastic collisions. Let us look briefly at tennis. Recall that in a collision, it is momentum and not force that is important. So, a heavier tennis racquet will have the advantage over a lighter one. This conclusion also holds true for other sports—a lightweight bat (such as a softball bat) cannot hit a hardball very far.\n\nThe location of the impact of the tennis ball on the racquet is also important, as is the part of the stroke during which the impact occurs. A smooth motion results in the maximizing of the velocity of the ball after impact and reduces sports injuries such as tennis elbow. A tennis player tries to hit the ball on the “sweet spot” on the racquet, where the vibration and impact are minimized and the ball is able to be given more velocity. Sports science and technologies also use physics concepts such as momentum and rotational motion and vibrations.\n\nTake-Home Experiment—Bouncing of Tennis Ball\n\n1. Find a racquet (a tennis, badminton, or other racquet will do). Place the racquet on the floor and stand on the handle. Drop a tennis ball on the strings from a measured height. Measure how high the ball bounces. Now ask a friend to hold the racquet firmly by the handle and drop a tennis ball from the same measured height above the racquet. Measure how high the ball bounces and observe what happens to your friend’s hand during the collision. Explain your observations and measurements.\n2. The coefficient of restitution $$(c)$$ is a measure of the elasticity of a collision between a ball and an object, and is defined as the ratio of the speeds after and before the collision. A perfectly elastic collision has a $$c$$ of 1. For a ball bouncing off the floor (or a racquet on the floor), $$c$$ can be shown to be $$c = (h/H)^{1/2}$$ where $$h$$ is the height to which the ball bounces and $$H$$ is the height from which the ball is dropped. Determine $$c$$ for the cases in Part 1 and for the case of a tennis ball bouncing off a concrete or wooden floor ($$c = 0.85$$ for new tennis balls used on a tennis court).\n\nExample $$\\PageIndex{2}$$: Calculating Final Velocity and Energy Release - Two Carts Collide\n\nIn the collision pictured in Figure $$\\PageIndex{3}$$, two carts collide inelastically. Cart 1 (denoted $$m_1$$ carries a spring which is initially compressed. During the collision, the spring releases its potential energy and converts it to internal kinetic energy. The mass of cart 1 and the spring is 0.350 kg, and the cart and the spring together have an initial velocity of $$-0.500 \\, m/s$$. After the collision, cart 1 is observed to recoil with a velocity of $$-4.00 \\, m/s$$.\n\n1. What is the final velocity of cart 2?\n2. How much energy was released by the spring (assuming all of it was converted into internal kinetic energy)?\n\nStrategy\n\nWe can use conservation of momentum to find the final velocity of cart 2, because $$F_{net} = 0$$ (the track is frictionless and the force of the spring is internal). Once this velocity is determined, we can compare the internal kinetic energy before and after the collision to see how much energy was released by the spring.\n\nSolution for (a)\n\nAs before, the equation for conservation of momentum in a two-object system is\n\n$m_1v_1 + m_2v_2 = m_1v'_1 + m_2v'_2. \\nonumber$\n\nThe only unknown in this equation is $$v'_2.$$ Solving for $$v'_2$$ and substituting known values into the previous equation fields\n\n\\begin{align*} v'_2 &= \\dfrac{m_1v_1 + m_2v_2 - m_1v'_1}{m_2} \\\\[5pt] &= \\dfrac{0.350 \\, kg)(2.00 \\, m/s) + (0.500 \\, kg)(-0.500 \\, m/s)}{0.500 \\, kg} - \\dfrac{(0.350 \\, kg)(-4.00 \\, m/s)}{0.500 \\, kg} \\\\[5pt] &= 3.70 \\, m/s.\\end{align*}\n\nSolution for (b)\n\nThe internal kinetic energy before the collision is\n\n\\begin{align*} KE_{int} &= \\dfrac{1}{2}m_1v_1^2 + \\dfrac{1}{2} m_2v_2^2 \\\\[5pt] &= \\dfrac{1}{2}(0.350 \\, kg)(2.00 \\, m/s)^2 + \\dfrac{1}{2}(0.500 \\, kg)(-0.500 \\, m/s)^2 \\\\[5pt] &= 0.763 \\, J. \\end{align*}\n\nAfter the collision, the internal kinetic energy is\n\n\\begin{align*} KE'_{int} &= \\dfrac{1}{2}m_1v_1^{'2} + \\dfrac{1}{2} m_2v_2^{'2} \\\\[5pt] &= \\dfrac{1}{2}(0.350 \\, kg)(-4.00 \\, m/s)^2 + \\dfrac{1}{2}(0.500 \\, kg)(0.370 \\, m/s)^2 \\\\[5pt] &= 6.22 \\, J. \\end{align*}\n\nThe change in internal kinetic energy is thus\n\n\\begin{align*} KE' - KE &= 6.22 \\, J - 0.763 \\, J \\\\[5pt] &= 5.46 \\, J. \\end{align*}\n\nDiscussion\n\nThe final velocity of cart 2 is large and positive, meaning that it is moving to the right after the collision. The internal kinetic energy in this collision increases by 5.46 J. That energy was released by the spring.\n\n## Summary\n\n• An inelastic collision is one in which the internal kinetic energy changes (it is not conserved).\n• A collision in which the objects stick together is sometimes called perfectly inelastic because it reduces internal kinetic energy more than does any other type of inelastic collision.\n• Sports science and technologies also use physics concepts such as momentum and rotational motion and vibrations.\n\n## Glossary\n\ninelastic collision\na collision in which internal kinetic energy is not conserved\nperfectly inelastic collision\na collision in which the colliding objects stick together" ]
[ null, "https://biz.libretexts.org/@api/deki/files/5084/girl-160172__340.png", null, "https://phys.libretexts.org/@api/deki/files/12577/imageedit_1_7213856719.jpg", null, "https://phys.libretexts.org/@api/deki/files/12578/imageedit_6_6268046528.png", null, "https://phys.libretexts.org/@api/deki/files/12579/imageedit_11_9666272580.png", null ]
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http://www.problemsphysics.com/mechanics/motion/disp_dist_solutions.html
[ "# Displacement and Distance: Solutions to Problems\n\nSolutions to the problems on displacement and distance of moving objects.\n\nSolutions to the problems on displacement and distance of moving objects.\n\n### Problem 1\n\nAn object moves from point A to point B to point C, then back to point B and then to point C along the line shown in the figure below.\na) Find the distance covered by the moving object.\nb) Find the magnitude and direction of the displacement of the object.", null, "Solution to Problem 1:\na) distance = AB + BC + CB + BC = 5 + 4 + 4 + 4 = 17 km\nb) The magnitude of the displacement is equal to the distance between the final point C and the initial point A = AC = 9 km\nThe direction of the displacement is the direction of the ray AB.\n\n### Problem 2\n\nAn object moves from point A to point C along the rectangle shown in the figure below.\na) Find the distance covered by the moving object.\nb) Find the magnitude of the displacement of the object.", null, "Solution to Problem 2:\na) distance = AB + BC = 5 + 3 = 8 km\nb) Initial point is A and the final point is C, hence the magnitude of the displacement is equal the diagonal AC of the rectangle and is calculated using Pythagora's theorem as follows\nAC\n2 = AB2 + BC2 = 52 + 32 = 25 + 9 = 34\nAC = √34 km\n\n### Problem 3\n\nAn object moves from point A to B to C to D and finally to A along the circle shown in the figure below.\na) Find the distance covered by the moving object.\nb) Find the magnitude and direction of the displacement of the object.", null, "Solution to Problem 3:\na) The object moves one complete rotation and therefore the distance d is equal to the circumference and is given by\nd = 2 Pi * radius = 6 Pi km\nb) Initial point is A and the final point is A, no change in position; hence the magnitude of the displacement is equal to zero\n\n### Problem 4\n\nAn object moves from point A to B to C to D along the circle shown in the figure below.\na) Find the distance covered by the moving object.\nb) Find the magnitude of the displacement of the object.", null, "Solution to Problem 4:\na) The object moves 3/4 of one rotation and therefore the distance d is equal to 3/4 of the circumference and is given by\nd = (3/4) (2 Pi * radius) = 4.5 Pi m\nb) Initial point is A and the final point is D, hence the magnitude of the displacement is equal to the distance AD which is calculated using Pythagora's theorem to the triangle AOD as shown in the figure below", null, "2 = AO2 + OD2 = 32 + 32 = 18\nmagnitude of displacement = AD = 3√2 m\n\n### Problem 5\n\nAn object moves along the grid through the points A, B, C, D, E, and F as shown below.\na) Find the distance covered by the moving object.\nb) Find the magnitude of the displacement of the object.", null, "Solution to Problem 5:\na) distance = AB + BC + CD + DE + EF = 3 + 1 + 1.5 + 0.5 + 0.5 = 6.5 km\nb) Initial point is A and the final point is F, hence the magnitude of the displacement is equal to the distance AF which is calculated by applying Pythagora's theorem to the triangle AHF as shown in the figure below", null, "AF\n2 = AH2 + HF2 = (0.5*4)2 + (0.5*3)2 = 4 + 2.25 = 6.25\nmagnitude of displacement = AF = 2.5 km" ]
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https://www.imath.kiev.ua/~sigma/2017/084/
[ "### Symmetry, Integrability and Geometry: Methods and Applications (SIGMA)\n\nSIGMA 13 (2017), 084, 21 pages      arXiv:1609.05270      https://doi.org/10.3842/SIGMA.2017.084\n\n### Realization of $U_q({\\mathfrak{sp}}_{2n})$ within the Differential Algebra on Quantum Symplectic Space\n\nJiao Zhang a and Naihong Hu b\na) Department of Mathematics, Shanghai University, Baoshan Campus, Shangda Road 99, Shanghai 200444, P.R. China\nb) Department of Mathematics, Shanghai Key Laboratory of Pure Mathematics and Mathematical Practice, East China Normal University, Minhang Campus, Dong Chuan Road 500, Shanghai 200241, P.R. China\n\nReceived April 18, 2017, in final form October 20, 2017; Published online October 27, 2017\n\nAbstract\nWe realize the Hopf algebra $U_q({\\mathfrak {sp}}_{2n})$ as an algebra of quantum differential operators on the quantum symplectic space $\\mathcal{X}(f_s;\\mathrm{R})$ and prove that $\\mathcal{X}(f_s;\\mathrm{R})$ is a $U_q({\\mathfrak{sp}}_{2n})$-module algebra whose irreducible summands are just its homogeneous subspaces. We give a coherence realization for all the positive root vectors under the actions of Lusztig's braid automorphisms of $U_q({\\mathfrak {sp}}_{2n})$.\n\nKey words: quantum symplectic group; quantum symplectic space; quantum differential operators; differential calculus; module algebra.\n\npdf (459 kb)   tex (24 kb)\n\nReferences\n\n1. Chari V., Xi N., Monomial bases of quantized enveloping algebras, in Recent Developments in Quantum Affine Algebras and Related Topics (Raleigh, NC, 1998), Contemp. Math., Vol. 248, Amer. Math. Soc., Providence, RI, 1999, 69-81, math.QA/9810167.\n2. Fiore G., Realization of $U_q({\\rm so}(N))$ within the differential algebra on ${\\bf R}^N_q$, Comm. Math. Phys. 169 (1995), 475-500, hep-th/9403033.\n3. Gu H., Hu N., Loewy filtration and quantum de Rham cohomology over quantum divided power algebra, J. Algebra 435 (2015), 1-32.\n4. Heckenberger I., Spin geometry on quantum groups via covariant differential calculi, Adv. Math. 175 (2003), 197-242, math.QA/0006226.\n5. Hu H., Hu N., Double-bosonization and Majid's conjecture, (I): Rank-inductions of $ABCD$, J. Math. Phys. 56 (2015), 111702, 16 pages, arXiv:1505.02612.\n6. Hu N., Quantum divided power algebra, $q$-derivatives, and some new quantum groups, J. Algebra 232 (2000), 507-540, arXiv:0902.2858.\n7. Hu N., Quantum group structure associated to the quantum affine space, Algebra Colloq. 11 (2004), 483-492.\n8. Jantzen J.C., Lectures on quantum groups, Graduate Studies in Mathematics, Vol. 6, Amer. Math. Soc., Providence, RI, 1996.\n9. Kassel C., Quantum groups, Graduate Texts in Mathematics, Vol. 155, Springer-Verlag, New York, 1995.\n10. Klimyk A., Schmüdgen K., Quantum groups and their representations, Texts and Monographs in Physics, Springer-Verlag, Berlin, 1997.\n11. Li Y., Hu N., The Green rings of the 2-rank Taft algebra and its two relatives twisted, J. Algebra 410 (2014), 1-35, arXiv:1305.1444.\n12. Lusztig G., Quantum deformations of certain simple modules over enveloping algebras, Adv. Math. 70 (1988), 237-249.\n13. Lusztig G., Introduction to quantum groups, Progress in Mathematics, Vol. 110, Birkhäuser Boston, Inc., Boston, MA, 1993.\n14. Majid S., Double-bosonization of braided groups and the construction of $U_q({\\mathfrak g})$, Math. Proc. Cambridge Philos. Soc. 125 (1999), 151-192, q-alg/9511001.\n15. Manin Yu.I., Quantum groups and noncommutative geometry, Université de Montréal, Centre de Recherches Mathématiques, Montreal, QC, 1988.\n16. Montgomery S., Smith S.P., Skew derivations and $U_q({\\rm sl}(2))$, Israel J. Math. 72 (1990), 158-166.\n17. Ogievetsky O., Differential operators on quantum spaces for ${\\rm GL}_q(n)$ and ${\\rm SO}_q(n)$, Lett. Math. Phys. 24 (1992), 245-255.\n18. Ogievetsky O., Schmidke W.B., Wess J., Zumino B., Six generator $q$-deformed Lorentz algebra, Lett. Math. Phys. 23 (1991), 233-240.\n19. Ogievetsky O., Schmidke W.B., Wess J., Zumino B., $q$-deformed Poincaré algebra, Comm. Math. Phys. 150 (1992), 495-518.\n20. Radford D.E., Towber J., Yetter-Drinfel'd categories associated to an arbitrary bialgebra, J. Pure Appl. Algebra 87 (1993), 259-279.\n21. Reshetikhin N.Yu., Takhtadzhyan L.A., Faddeev L.D., Quantization of Lie groups and Lie algebras, Leningrad Math. J. 1 (1990), 178-206.\n22. Schauenburg P., Hopf modules and Yetter-Drinfel'd modules, J. Algebra 169 (1994), 874-890.\n23. Wess J., Zumino B., Covariant differential calculus on the quantum hyperplane, Nuclear Phys. B Proc. Suppl. 18 (1990), 302-312.\n24. Woronowicz S.L., Differential calculus on compact matrix pseudogroups (quantum groups), Comm. Math. Phys. 122 (1989), 125-170.\n25. Yetter D.N., Quantum groups and representations of monoidal categories, Math. Proc. Cambridge Philos. Soc. 108 (1990), 261-290." ]
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https://www.convertunits.com/from/kilojoules/to/joule
[ "## ››Convert kilojoule to joule\n\n kilojoules joule\n\nHow many kilojoules in 1 joule? The answer is 0.001.\nWe assume you are converting between kilojoule and joule.\nYou can view more details on each measurement unit:\nkilojoules or joule\nThe SI derived unit for energy is the joule.\n1 kilojoules is equal to 1000 joule.\nNote that rounding errors may occur, so always check the results.\nUse this page to learn how to convert between kilojoules and joules.\nType in your own numbers in the form to convert the units!\n\n## ››Quick conversion chart of kilojoules to joule\n\n1 kilojoules to joule = 1000 joule\n\n2 kilojoules to joule = 2000 joule\n\n3 kilojoules to joule = 3000 joule\n\n4 kilojoules to joule = 4000 joule\n\n5 kilojoules to joule = 5000 joule\n\n6 kilojoules to joule = 6000 joule\n\n7 kilojoules to joule = 7000 joule\n\n8 kilojoules to joule = 8000 joule\n\n9 kilojoules to joule = 9000 joule\n\n10 kilojoules to joule = 10000 joule\n\n## ››Want other units?\n\nYou can do the reverse unit conversion from joule to kilojoules, or enter any two units below:\n\n## Enter two units to convert\n\n From: To:\n\n## ››Definition: Kilojoule\n\nThe SI prefix \"kilo\" represents a factor of 103, or in exponential notation, 1E3.\n\nSo 1 kilojoule = 103 joules.\n\nThe definition of a joule is as follows:\n\nThe joule (symbol J, also called newton meter, watt second, or coulomb volt) is the SI unit of energy and work. The unit is pronounced to rhyme with \"tool\", and is named in honor of the physicist James Prescott Joule (1818-1889).\n\n## ››Definition: Joule\n\nThe joule (symbol J, also called newton meter, watt second, or coulomb volt) is the SI unit of energy and work. The unit is pronounced to rhyme with \"tool\", and is named in honor of the physicist James Prescott Joule (1818-1889).\n\n## ››Metric conversions and more\n\nConvertUnits.com provides an online conversion calculator for all types of measurement units. You can find metric conversion tables for SI units, as well as English units, currency, and other data. Type in unit symbols, abbreviations, or full names for units of length, area, mass, pressure, and other types. Examples include mm, inch, 100 kg, US fluid ounce, 6'3\", 10 stone 4, cubic cm, metres squared, grams, moles, feet per second, and many more!" ]
[ null ]
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https://doc.cgal.org/5.4.4/AABB_tree/classCGAL_1_1AABB__traits.html
[ "", null, "CGAL 5.4.4 - 3D Fast Intersection and Distance Computation (AABB Tree)\nCGAL::AABB_traits< GeomTraits, AABBPrimitive, BboxMap > Class Template Reference\n\n#include <CGAL/AABB_traits.h>\n\n## Definition\n\n### template<typename GeomTraits, typename AABBPrimitive, typename BboxMap = Default> class CGAL::AABB_traits< GeomTraits, AABBPrimitive, BboxMap >\n\nThis traits class handles any type of 3D geometric primitives provided that the proper intersection tests and constructions are implemented.\n\nIt handles points, rays, lines and segments as query types for intersection detection and computations, and it handles points as query type for distance queries.\n\nIs Model Of:\nTemplate Parameters\n GeomTraits must be a model of the concept AABBGeomTraits, and provide the geometric types as well as the intersection tests and computations. Primitive provide the type of primitives stored in the AABB_tree. It is a model of the concept AABBPrimitive or AABBPrimitiveWithSharedData. BboxMap must be a model of ReadablePropertyMap that has as key type a primitive id, and as value type a Bounding_box. If the type is Default the Datum must have the member function bbox() that returns the bounding box of the primitive.\n\nIf the argument GeomTraits is a model of the concept AABBRayIntersectionGeomTraits, this class is also a model of AABBRayIntersectionTraits.\n\nAABBTraits\nAABB_tree\nAABBPrimitive\nAABBPrimitiveWithSharedData\nExamples:\nAABB_tree/AABB_custom_example.cpp, AABB_tree/AABB_custom_indexed_triangle_set_array_example.cpp, AABB_tree/AABB_custom_indexed_triangle_set_example.cpp, AABB_tree/AABB_custom_triangle_soup_example.cpp, AABB_tree/AABB_face_graph_triangle_example.cpp, AABB_tree/AABB_halfedge_graph_edge_example.cpp, AABB_tree/AABB_insertion_example.cpp, AABB_tree/AABB_polyhedron_edge_example.cpp, AABB_tree/AABB_polyhedron_facet_distance_example.cpp, AABB_tree/AABB_polyhedron_facet_intersection_example.cpp, AABB_tree/AABB_ray_shooting_example.cpp, AABB_tree/AABB_segment_3_example.cpp, and AABB_tree/AABB_triangle_3_example.cpp.\n\n## Classes\n\nclass  Do_intersect\nFunction object using GeomTraits::Do_intersect. More...\n\nstruct  Intersection_and_primitive_id\nIntersection_and_primitive_id<Query>::Type::first_type is found according to the result type of GeomTraits::Intersect_3::operator(). More...\n\n## Public Member Functions\n\nAABB_traits ()\nDefault constructor.\n\n## Types\n\ntypedef GeomTraits::Point_3 Point_3\nPoint query type.\n\ntypedef GeomTraits::Iso_cuboid_3 Iso_cuboid_3\nadditionnal types for the search tree, required by the RangeSearchTraits concept More...\n\ntypedef CGAL::Bbox_3 Bounding_box\nBounding box type.\n\n## ◆ Iso_cuboid_3\n\ntemplate<typename GeomTraits , typename AABBPrimitive , typename BboxMap = Default>\n typedef GeomTraits::Iso_cuboid_3 CGAL::AABB_traits< GeomTraits, AABBPrimitive, BboxMap >::Iso_cuboid_3\n\nadditionnal types for the search tree, required by the RangeSearchTraits concept" ]
[ null, "https://doc.cgal.org/5.4.4/Manual/search/mag_sel.png", null ]
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https://www.go2live.cn/nocate/%E9%93%BE%E8%A1%A8%E3%80%81%E6%A0%88%E3%80%81%E9%98%9F%E5%88%97%E3%80%81kmp%E7%9B%B8%E5%85%B3%E7%9F%A5%E8%AF%86%E7%82%B9.html
[ "链表、栈、队列、KMP相关知识点\n\n数组模拟单链表:\n\n1. 比使用结构体 或者类来说 速度更快\n2. 代码简洁\n3. 算法题:空间换时间", null, "import java.util.*;\n\npublic class Main{\nstatic int N = 100010;\nstatic int[] e = new int[N];\nstatic int[] next = new int[N];\n\n// 初始化\nstatic void init(){\nidx = 0;\n}\n// 向链表头插入一个数\ne[idx] = x;\n}\n\n// 向k位置插入x\nstatic void insert(int k,int x){\ne[idx] = x;\nnext[idx] = next[k];\nnext[k] = idx++;\n}\n\n// 删除k位置的数\nstatic void delete(int k){\nnext[k] = next[next[k]];\n}\n// 主函数\npublic static void main(String[] args){\nScanner sc = new Scanner(System.in);\ninit();\nint n = sc.nextInt();\nwhile(n-- != 0){\nString s = sc.next();\nif(s.equals(\"H\")){\nint x = sc.nextInt();\n}else if(s.equals(\"D\")){\nint k = sc.nextInt();\nelse delete(k-1);\n}else if(s.equals(\"I\")){\nint k = sc.nextInt();\nint x = sc.nextInt();\ninsert(k-1,x);\n}\n\n}\nSystem.out.print(e[i]+\" \");\n}\n\n}\n\n}\n\n栈:\n\n// tt表示栈顶\nint stk[N], tt = 0;\n\n// 向栈顶插入一个数\nstk[ ++ tt] = x;\n\n// 从栈顶弹出一个数\ntt -- ;\n\n// 栈顶的值\nstk[tt];\n\n// 判断栈是否为空\nif (tt > 0)\n{\n\n}\n\n队列:\n\n1. 普通队列:\n// hh 表示队头,tt表示队尾\nint q[N], hh = 0, tt = -1;\n\n// 向队尾插入一个数\nq[ ++ tt] = x;\n\n// 从队头弹出一个数\nhh ++ ;\n\n// 队头的值\nq[hh];\n\n// 判断队列是否为空\nif (hh <= tt)\n{\n\n}\n2. 循环队列\n// hh 表示队头,tt表示队尾的后一个位置\nint q[N], hh = 0, tt = 0;\n\n// 向队尾插入一个数\nq[tt ++ ] = x;\nif (tt == N) tt = 0;\n\n// 从队头弹出一个数\nhh ++ ;\nif (hh == N) hh = 0;\n\n// 队头的值\nq[hh];\n\n// 判断队列是否为空\nif (hh != tt)\n{\n\n}\n\n单调栈:\n\nimport java.util.*;\nimport java.io.*;\n\npublic class Main{\nstatic int N = 100010;\nstatic int[] q = new int[N];\n\npublic static void main(String[] args){\nScanner sc = new Scanner(System.in);\nint n = sc.nextInt();\nfor(int i = 0;i < n; i++){\nq[i] = sc.nextInt();\n}\n\nfor(int i = 0;i =0;j--){\nif(q[i]>q[j]){\nindex = j;\nSystem.out.print(q[j]+\" \");\nbreak;\n}\n}\nif(q[i]<=q[index]){\nSystem.out.print(\"-1 \");\n}\n}\n\n}\n\n}\n\nAC代码:\n\nimport java.util.*;\n\npublic class Main{\nstatic int N = 100010;\nstatic int[] stk = new int[N];\nstatic int tt=0;\npublic static void main(String[] args){\nScanner sc = new Scanner(System.in);\n\nint n = sc.nextInt();\n\nfor(int i = 0; i = x) tt--;\n// 上面结束以后是栈顶元素小于输入的元素; stk[tt] < x; 这样就保证了x元素在左边并且小于本身的元素是栈顶元素;\n//如果栈中有元素,输出栈顶元素\nif(tt != 0) System.out.print(stk[tt]+\" \");\n// 否则栈中没有元素\nelse System.out.print(\"-1\"+\" \");\n// 将元素加入栈中\nstk[++tt] = x;\n}\n\n}\n\n}\n\n单调队列:\n\n1. 判断队头出没出窗口 if true–>hh++;\n2. 求最小值,保持队列单调上升,判断队尾元素tt与当前元素a[i]的大小,若tt >=a[i],剔除队尾元素;\n3. 求最大值,保持队列单调下降,判断队尾元素tt与当前元素a[i]的大小,若tt <=a[i],剔除队尾元素;\n4. 将当前元素下标加入队尾;\n5. 如果满足条件输出\n\nimport java.io.*;\npublic class Main{\nstatic int N = 1000010;\n// 原数组\nstatic int[] a = new int[N];\n// 队列\nstatic int[] q = new int[N];\n\nstatic int hh=0,tt=-1;\npublic static void main(String[] args) throws IOException{\nBufferedWriter bw = new BufferedWriter(new OutputStreamWriter(System.out));\n\n// a数组中由n个元素\nint n = Integer.parseInt(num);\n\n// 滑动窗口的大小\nint k = Integer.parseInt(num);\n\n// 初始化a数组\nfor(int i = 0; i < n; i++) a[i] = Integer.parseInt(nums[i]);\n\n// 查找最小值\nfor(int i = 0; i < n; i++){\n\n// 判断当前窗口是否大于滑动窗口的大小\nif(hh k) hh++;\n\n// 判断如果有元素并且队尾元素大于入队元素,则舍弃队尾元素来保证单调上升的队列;\nwhile(hh = a[i]) tt--;\n\n//向队列加入元素下标\nq[++tt] = i;\n\n// 如果下标大于等于窗口大小,那么输出队头元素\nif(i+1 >=k) bw.write(a[q[hh]]+\" \");\n}\nbw.write(\"\\n\");\nhh = 0;\ntt = -1;\n// 查找最大值\nfor(int i = 0; i < n; i++){\n\n// 判断当前窗口是否大于滑动窗口的大小\nif(hh k) hh++;\n\n// 判断如果有元素并且队尾元素小于入队元素,则舍弃队尾元素来保证单调上升的队列;\nwhile(hh <= tt && a[q[tt]] =k) bw.write(a[q[hh]]+\" \");\n}\nbw.flush();\nbr.close();\nbw.close();\n}\n\n}\n\nKMP算法:\n\nKMP定义:是取自三个发明人的首字母组成的;\n\nkmp实现的方式:求next[]、以及匹配字符串\n\n前缀与后缀的概念:(很重要)\n\nababa\n\na,ab,aba,abab\n\nbaba,aba,ba,a", null, "", null, "$\\because 匹配失败使得p1移动到p2位置\\\\ \\therefore s-1 = p2-2.2\\\\ \\because p2是p1的平移\\\\ \\therefore p1 = p2\\\\ \\therefore p1-2.1 = p2-2.2\\\\ 又\\because s = p1\\\\ \\therefore s-1 = p1-3\\\\ 由上可知:\\\\ p2-2.2 = p1-3;\\\\ p1-2.1 = p1-3;$\n\n// kmp匹配过程,遍历s模板每个元素\nfor(int i =1, j =0; i<=m; i++ ){\n// 如果j回退到0或者i位置元素与j+1位置的元素不相同,那么执行回退操作,\n//j退回next[i]处,即前缀与后缀相同的区间最后元素位置\nwhile(j != 0 && s[i] != p[j+1]) j = next[j];\nif(s[i] == p[j+1]) j++;\n// 如果匹配成功\nif(j ==n){\n// 输出匹配元素在s模板中的起始位置\nbw.write(i-n+\" \");\n// 继续匹配;\nj = next[j];\n}\n\nnext[]数组(难点)\n\np a b c a b\n\nnext[ ] 0 0 0 1 2", null, "// 实现next数组(找到前缀与后缀相同的最大元素长度)\nint[] next = new int[n+1];\n//next = 0;\nfor(int i = 2,j =0; i<= n;i++){\n// 如果j没有回退到0并且i位置元素与j+1位置的元素不相同,那么执行回退操作,\n//j退回next[i]处,即前缀与后缀相同的区间最后元素位置\nwhile(j !=0 && p[i] != p[j+1]) j = next[j];\n// 如果i与j+1相同,那么移动j向后匹配\nif(p[i] == p[j+1]) j++;\n// p[1,j] = p[i-j+1,i];前缀与后缀相同;i表示终点\nnext[i] = j;\n}\n\n完整代码:\n\nimport java.util.*;\nimport java.io.*;\n\n//下标为什么从1开始,简化代码的复杂度;\npublic class Main{\n\npublic static void main(String[] args) throws IOException{\nBufferedWriter bw = new BufferedWriter(new OutputStreamWriter(System.out));\n\n// 对模式串p进行操作\nchar[] p = new char[n+1];\n// 将输入的字符保存到缓冲区\nfor(int i = 1; i <=n;i++){\n// 取出的字符串下标-1;\np[i] = pstr.charAt(i-1);\n}\n\n// 对模板串进行操作\nchar[] s = new char[m+1];\nfor(int i = 1;i <= m;i++){\ns[i] = sstr.charAt(i-1);\n}\n\n// 实现next数组\nint[] next = new int[n+1];\n//next = 0;\nfor(int i = 2,j =0; i<= n;i++){\n// 如果j回退到0或者i位置元素与j+1位置的元素不相同,那么执行回退操作,\n//j退回next[i]处,即前缀与后缀相同的区间最后元素位置\nwhile(j !=0 && p[i] != p[j+1]) j = next[j];\n// 如果i与j+1相同,那么移动j向后匹配\nif(p[i] == p[j+1]) j++;\n// p[1,j] = p[i-j+1,i];前缀与后缀相同;i表示终点\nnext[i] = j;\n}\n\n// kmp匹配过程,遍历s模板每个元素\nfor(int i =1, j =0; i<=m; i++ ){\n// 如果j回退到0或者i位置元素与j+1位置的元素不相同,那么执行回退操作,\n//j退回next[i]处,即前缀与后缀相同的区间最后元素位置\nwhile(j != 0 && s[i] != p[j+1]) j = next[j];\nif(s[i] == p[j+1]) j++;\n// 如果匹配成功\nif(j ==n){\n// 输出匹配元素在s模板中的起始位置\nbw.write(i-n+\" \");\n// 继续匹配;\nj = next[j];\n}\n\n}\nbw.flush();\nbr.close();\nbw.close();\n}\n\n}" ]
[ null, "https://img0.tuicool.com/ZfeiUbF.png!web", null, "https://img0.tuicool.com/aQnIfyJ.gif!web", null, "https://img0.tuicool.com/zQJJfyq.png!web", null, "https://img2.tuicool.com/vIzIJv.png!web", null ]
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http://eel.is/c++draft/unord.map.overview
[ "# 22 Containers library [containers]\n\n## 22.5 Unordered associative containers [unord]\n\n### 22.5.4 Class template unordered_­map[unord.map]\n\n#### 22.5.4.1 Overview [unord.map.overview]\n\nAn unordered_­map is an unordered associative container that supports unique keys (an unordered_­map contains at most one of each key value) and that associates values of another type mapped_­type with the keys.\nThe unordered_­map class supports forward iterators.\nAn unordered_­map meets all of the requirements of a container, of an unordered associative container, and of an allocator-aware container (Table 78).\nIt provides the operations described in the preceding requirements table for unique keys; that is, an unordered_­map supports the a_­uniq operations in that table, not the a_­eq operations.\nFor an unordered_­map<Key, T> the key type is Key, the mapped type is T, and the value type is pair<const Key, T>.\nSubclause [unord.map] only describes operations on unordered_­map that are not described in one of the requirement tables, or for which there is additional semantic information.\nnamespace std { template<class Key, class T, class Hash = hash<Key>, class Pred = equal_to<Key>, class Allocator = allocator<pair<const Key, T>>> class unordered_map { public: // types using key_type = Key; using mapped_type = T; using value_type = pair<const Key, T>; using hasher = Hash; using key_equal = Pred; using allocator_type = Allocator; using pointer = typename allocator_traits<Allocator>::pointer; using const_pointer = typename allocator_traits<Allocator>::const_pointer; using reference = value_type&; using const_reference = const value_type&; using size_type = implementation-defined; // see [container.requirements] using difference_type = implementation-defined; // see [container.requirements] using iterator = implementation-defined; // see [container.requirements] using const_iterator = implementation-defined; // see [container.requirements] using local_iterator = implementation-defined; // see [container.requirements] using const_local_iterator = implementation-defined; // see [container.requirements] using node_type = unspecified; using insert_return_type = insert-return-type<iterator, node_type>; // [unord.map.cnstr], construct/copy/destroy unordered_map(); explicit unordered_map(size_type n, const hasher& hf = hasher(), const key_equal& eql = key_equal(), const allocator_type& a = allocator_type()); template<class InputIterator> unordered_map(InputIterator f, InputIterator l, size_type n = see below, const hasher& hf = hasher(), const key_equal& eql = key_equal(), const allocator_type& a = allocator_type()); unordered_map(const unordered_map&); unordered_map(unordered_map&&); explicit unordered_map(const Allocator&); unordered_map(const unordered_map&, const Allocator&); unordered_map(unordered_map&&, const Allocator&); unordered_map(initializer_list<value_type> il, size_type n = see below, const hasher& hf = hasher(), const key_equal& eql = key_equal(), const allocator_type& a = allocator_type()); unordered_map(size_type n, const allocator_type& a) : unordered_map(n, hasher(), key_equal(), a) { } unordered_map(size_type n, const hasher& hf, const allocator_type& a) : unordered_map(n, hf, key_equal(), a) { } template<class InputIterator> unordered_map(InputIterator f, InputIterator l, size_type n, const allocator_type& a) : unordered_map(f, l, n, hasher(), key_equal(), a) { } template<class InputIterator> unordered_map(InputIterator f, InputIterator l, size_type n, const hasher& hf, const allocator_type& a) : unordered_map(f, l, n, hf, key_equal(), a) { } unordered_map(initializer_list<value_type> il, size_type n, const allocator_type& a) : unordered_map(il, n, hasher(), key_equal(), a) { } unordered_map(initializer_list<value_type> il, size_type n, const hasher& hf, const allocator_type& a) : unordered_map(il, n, hf, key_equal(), a) { } ~unordered_map(); unordered_map& operator=(const unordered_map&); unordered_map& operator=(unordered_map&&) noexcept(allocator_traits<Allocator>::is_always_equal::value && is_nothrow_move_assignable_v<Hash> && is_nothrow_move_assignable_v<Pred>); unordered_map& operator=(initializer_list<value_type>); allocator_type get_allocator() const noexcept; // iterators iterator begin() noexcept; const_iterator begin() const noexcept; iterator end() noexcept; const_iterator end() const noexcept; const_iterator cbegin() const noexcept; const_iterator cend() const noexcept; // capacity [[nodiscard]] bool empty() const noexcept; size_type size() const noexcept; size_type max_size() const noexcept; // [unord.map.modifiers], modifiers template<class... Args> pair<iterator, bool> emplace(Args&&... args); template<class... Args> iterator emplace_hint(const_iterator position, Args&&... args); pair<iterator, bool> insert(const value_type& obj); pair<iterator, bool> insert(value_type&& obj); template<class P> pair<iterator, bool> insert(P&& obj); iterator insert(const_iterator hint, const value_type& obj); iterator insert(const_iterator hint, value_type&& obj); template<class P> iterator insert(const_iterator hint, P&& obj); template<class InputIterator> void insert(InputIterator first, InputIterator last); void insert(initializer_list<value_type>); node_type extract(const_iterator position); node_type extract(const key_type& x); insert_return_type insert(node_type&& nh); iterator insert(const_iterator hint, node_type&& nh); template<class... Args> pair<iterator, bool> try_emplace(const key_type& k, Args&&... args); template<class... Args> pair<iterator, bool> try_emplace(key_type&& k, Args&&... args); template<class... Args> iterator try_emplace(const_iterator hint, const key_type& k, Args&&... args); template<class... Args> iterator try_emplace(const_iterator hint, key_type&& k, Args&&... args); template<class M> pair<iterator, bool> insert_or_assign(const key_type& k, M&& obj); template<class M> pair<iterator, bool> insert_or_assign(key_type&& k, M&& obj); template<class M> iterator insert_or_assign(const_iterator hint, const key_type& k, M&& obj); template<class M> iterator insert_or_assign(const_iterator hint, key_type&& k, M&& obj); iterator erase(iterator position); iterator erase(const_iterator position); size_type erase(const key_type& k); iterator erase(const_iterator first, const_iterator last); void swap(unordered_map&) noexcept(allocator_traits<Allocator>::is_always_equal::value && is_nothrow_swappable_v<Hash> && is_nothrow_swappable_v<Pred>); void clear() noexcept; template<class H2, class P2> void merge(unordered_map<Key, T, H2, P2, Allocator>& source); template<class H2, class P2> void merge(unordered_map<Key, T, H2, P2, Allocator>&& source); template<class H2, class P2> void merge(unordered_multimap<Key, T, H2, P2, Allocator>& source); template<class H2, class P2> void merge(unordered_multimap<Key, T, H2, P2, Allocator>&& source); // observers hasher hash_function() const; key_equal key_eq() const; // map operations iterator find(const key_type& k); const_iterator find(const key_type& k) const; template<class K> iterator find(const K& k); template<class K> const_iterator find(const K& k) const; template<class K> size_type count(const key_type& k) const; template<class K> size_type count(const K& k) const; bool contains(const key_type& k) const; template<class K> bool contains(const K& k) const; pair<iterator, iterator> equal_range(const key_type& k); pair<const_iterator, const_iterator> equal_range(const key_type& k) const; template<class K> pair<iterator, iterator> equal_range(const K& k); template<class K> pair<const_iterator, const_iterator> equal_range(const K& k) const; // [unord.map.elem], element access mapped_type& operator[](const key_type& k); mapped_type& operator[](key_type&& k); mapped_type& at(const key_type& k); const mapped_type& at(const key_type& k) const; // bucket interface size_type bucket_count() const noexcept; size_type max_bucket_count() const noexcept; size_type bucket_size(size_type n) const; size_type bucket(const key_type& k) const; local_iterator begin(size_type n); const_local_iterator begin(size_type n) const; local_iterator end(size_type n); const_local_iterator end(size_type n) const; const_local_iterator cbegin(size_type n) const; const_local_iterator cend(size_type n) const; // hash policy float load_factor() const noexcept; float max_load_factor() const noexcept; void max_load_factor(float z); void rehash(size_type n); void reserve(size_type n); }; template<class InputIterator, class Hash = hash<iter-key-type<InputIterator>>, class Pred = equal_to<iter-key-type<InputIterator>>, class Allocator = allocator<iter-to-alloc-type<InputIterator>>> unordered_map(InputIterator, InputIterator, typename see below::size_type = see below, Hash = Hash(), Pred = Pred(), Allocator = Allocator()) -> unordered_map<iter-key-type<InputIterator>, iter-mapped-type<InputIterator>, Hash, Pred, Allocator>; template<class Key, class T, class Hash = hash<Key>, class Pred = equal_to<Key>, class Allocator = allocator<pair<const Key, T>>> unordered_map(initializer_list<pair<Key, T>>, typename see below::size_type = see below, Hash = Hash(), Pred = Pred(), Allocator = Allocator()) -> unordered_map<Key, T, Hash, Pred, Allocator>; template<class InputIterator, class Allocator> unordered_map(InputIterator, InputIterator, typename see below::size_type, Allocator) -> unordered_map<iter-key-type<InputIterator>, iter-mapped-type<InputIterator>, hash<iter-key-type<InputIterator>>, equal_to<iter-key-type<InputIterator>>, Allocator>; template<class InputIterator, class Allocator> unordered_map(InputIterator, InputIterator, Allocator) -> unordered_map<iter-key-type<InputIterator>, iter-mapped-type<InputIterator>, hash<iter-key-type<InputIterator>>, equal_to<iter-key-type<InputIterator>>, Allocator>; template<class InputIterator, class Hash, class Allocator> unordered_map(InputIterator, InputIterator, typename see below::size_type, Hash, Allocator) -> unordered_map<iter-key-type<InputIterator>, iter-mapped-type<InputIterator>, Hash, equal_to<iter-key-type<InputIterator>>, Allocator>; template<class Key, class T, class Allocator> unordered_map(initializer_list<pair<Key, T>>, typename see below::size_type, Allocator) -> unordered_map<Key, T, hash<Key>, equal_to<Key>, Allocator>; template<class Key, class T, class Allocator> unordered_map(initializer_list<pair<Key, T>>, Allocator) -> unordered_map<Key, T, hash<Key>, equal_to<Key>, Allocator>; template<class Key, class T, class Hash, class Allocator> unordered_map(initializer_list<pair<Key, T>>, typename see below::size_type, Hash, Allocator) -> unordered_map<Key, T, Hash, equal_to<Key>, Allocator>; }\nA size_­type parameter type in an unordered_­map deduction guide refers to the size_­type member type of the type deduced by the deduction guide." ]
[ null ]
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https://planetmath.org/bornologicalspace
[ "# bornological space\n\nA bornivore is a set which absorbs all bounded sets. That is, $G$ is a bornivore if given any bounded set $B$, there exists a $\\delta>0$ such that $\\epsilon B\\subset G$ for $0\\leq\\epsilon<\\delta$.\n\n## References\n\n• 1\n• 2 H.H. Schaefer, M. P. Wolff, Topological Vector Spaces, 2nd ed. 1999, Springer-Verlag.\nTitle bornological space BornologicalSpace 2013-03-22 15:59:09 2013-03-22 15:59:09 Mathprof (13753) Mathprof (13753) 8 Mathprof (13753) Definition msc 46A08 bornivore" ]
[ null ]
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https://en.m.wikibooks.org/wiki/QBasic/Appendix
[ "# QBasic/Appendix\n\n## Commands\n\n#### ABS()\n\n```N = ABS(expression returning a numerical value)\n```\n\nReturns the 'absolute' value of the expression, turning a negative to a positive (e.g. -4 to 4)\n\n`PRINT ABS(54.345) 'This will print the value ABS now as it is (54.345)`\n`PRINT ABS(-43) 'This will print the value as (43)`\n\n.\n\n#### ACCESS\n\n`OPEN \"file.txt\" FOR APPEND ACCESS WRITE`\n\nThis sets the access of a file that has been declared into the program. There are three settings that the programmer can set. These are:\n\n```READ - Sets up the file to be read only, no writing.\nWRITE - Writes only to the file. Cannot be read.\nREAD WRITE - Sets the file to both of the settings above.\n```\n\nFOR APPEND opens the file as a text file and sets the file pointer to the end of the file, thereby appending new output to the file. If omitted, FOR RANDOM is the default, which would open the file as a sequence of fixed-length records and set the file pointer to the start of the file. New data would overwrite existing data in the file without warning.\n\n#### ASC(\"C\")\n\n`PRINT ASC(\"t\") 'Will print 116`\n\nPrints the ASCII code number of the character found within the brackets. If the programmer put in a string into brackets, only the first character of the string will be shown.\n\n#### ATN()\n\n```ATN(expression)\n```\n\nPart of the inbuilt trigonometry functions. An expression that evaluates to a numeric vale is converted to it's Arc-Tangent.\n\n```angle = ATN( B )\nangle2 = ATN( 23 / 34 )\n```\n\n#### BEEP\n\n```BEEP\n```\n\nThe BIOS on the motherboard is instructed to emit a \"Beep\" sound from the PC 'speaker'. See also SOUND and PLAY.\n\nAn older, outdated alternative is to use : PRINT CHR\\$(07)\n\nThis was replaced later by the BEEP command.\n\n```BLOAD file_path\\$, memory_offset%\n```\n\nLoads a file saved with BSAVE into memory.\n\n• file_path\\$ is the file location\n• memory_offset is an integer specifying the offset in memory to load the file to\n\nThe starting memory address is determined by the offset and the most recent call to the DEF SEG statement\n\n#### BSAVE\n\n```BSAVE file_path\\$, memory_offset%, length&\n```\n\nSaves the contents of an area in memory to a file that can be loaded with BLOAD\n\n• file_path\\$ is the file location\n• memory_offset is an integer specifying the offset in memory to save\n• length the number of bytes to copy\n\nThe starting memory address is determined by the offset and the most recent call to the DEF SEG statement\n\n#### CALL ABSOLUTE\n\n```CALL ABSOLUTE([argument%, ...,] address%)\n```\n\nPushes the provided arguments, which must be INTEGER, from left to right on the stack and then does a far call to the assembly language routine located at address. The code segment to use is set using DEF SEG. Normally QBasic will push the address of arguments, but if an argument is preceded by BYVAL the value of the argument will be pushed.\n\nNote that because QBasic pushes the arguments from left to right, if you provide three arguments for example the stack will look like this:\n\n```SS:SP Return IP\n+0002 Return CS\n+0004 Argument 3\n+0006 Argument 2\n+0008 Argument 1\n```\n\nExample:\n\n```'POP CX\n'POP DX\n'POP BX\n'POP AX\n'MOV [BX], AX\n'PUSH DX\n'PUSH CX\n'RETF\nA\\$ = \"YZ[Xë•RQ╦\" 'Codepage: 437\nI% = 42: J% = 0\nIF VARSEG(A\\$) <> VARSEG(J%) THEN STOP 'Both A\\$ and J% are stored in DGROUP.\nDEF SEG = VARSEG(A\\$)\nCALL ABSOLUTE(BYVAL I%, J%, PEEK(VARPTR(A\\$) + 3) * 256& OR PEEK(VARPTR(A\\$) + 2))\nPRINT J%\n```\n\nQBasic doesn't set BP to the stack frame of your machine language routine, but leaves it pointing to the stack frame of the calling QBasic procedure, which looks like this:\n\n```-???? Variable 3\n-???? Variable 2\n-???? Variable 1\n-???? Return value (only if the procedure is a FUNCTION, absent if a SUB)\n-0002 BP of calling procedure\nSS:BP BP of calling procedure (yes, again)\n+0002 Six bytes referring back to the calling procedure to use when executing END/EXIT SUB/FUNCTION\n+0008 Argument 3\n+000A Argument 2\n+000C Argument 1\n```\n\nThe offsets indicated with -/+???? will depend on the sizes and presence of variables, return value and arguments. For example, an INTEGER variable will take two bytes, but a LONG variable four. As one might expect considering how QBasic passes arguments, variables are stored in reverse order of declaration. Contrary to when calling a machine language routine, the arguments here will always be addresses. For arguments passed by value, space is allocated on the stack and the address of that space is passed to the procedure.\n\n#### CASE\n\n```SELECT CASE expression\nCASE test1[, test2, ...]: statements\n[CASE test4[, test5, ...]: statements]\n[CASE ELSE: statements]\nEND SELECT\n```\n\nExecutes the statements after the first CASE statement where a test matches expression. The tests can be of the following forms:\n\n```expression\nexpression1 TO expression2\nIS {<|<=|=|<>|>=|>} expression\n```\n\nThis is an example of a program with no CASE commands that assigns different paths to values:\n\n```PRINT \"1. Print 'path'\"\nPRINT \"2. Print 'hello'\"\nPRINT \"3. Quit\"\nINPUT \"Enter a choice: \"; a\\$\nIF a\\$ = \"1\" THEN PRINT \"path\": RUN\nIF a\\$ = \"2\" THEN PRINT \"hello\": RUN\nIF a\\$ <> \"3\" THEN PRINT \"That is not a valid choice.\": RUN\n```\n\nThis is what a program looks like with the CASE command:\n\n```PRINT \"1. Print 'path'\"\nPRINT \"2. Print 'hello'\"\nPRINT \"3. Quit\"\nINPUT \"Enter a choice: \"; a\\$\nSELECT CASE a\\$\nCASE \"1\": PRINT \"path\": RUN\nCASE \"2\": PRINT \"hello\": RUN\nCASE IS <> \"3\": PRINT \"That is not a valid choice.\": RUN\nEND SELECT\n```\n\n#### CHAIN\n\n```CHAIN filename\n```\n\nThis chains execution to another QBasic program. Values may be passed to the other program by using the COMMON statement before the CHAIN statement. Note that execution doesn't return to the first program unless the second program uses CHAIN to transfer execution back to the first.\n\n#### CHDIR\n\n```CHDIR directoryname\n```\n\nThis is used for setting the working directory, also known as the current directory. The directory name is declared exactly like in DOS and long file names aren't supported. For example:\n\n```CHDIR \"C:\\DOS\"\n```\n\nNote that this doesn't change the current drive and every drive has its own working directory. You can set the current drive like this:\n\n```SHELL \"C:\"\n```\n\n#### CHR\\$()\n\nThis returns the string character symbol of an ASCII code value.\n\n```name\\$ = CHR\\$([ascii character code])\n```\n\nOften used to 'load' characters into string variables when that character cannot be typed (e.g. the Esc key or the F{n} (Function Keys) or characters that would be 'recognised' and acted upon by the QBASIC interpreter. The following four characters cannot occur in a QBasic string literal:\n\n• 0 Null: all characters up to the end of the line will get deleted, including this one.\n• 10 Line Feed: signals the end of the line.\n• 13 Carriage Return: this character will get deleted.\n• 34 Quotation Mark: signals the end of the string literal.\n\nHere is a list of some character codes :-\n\n```07 Beep (same as BEEP)\n08 Backspace\n09 Tab\n27 Esc\n72 Up Arrow\n75 Left Arrow\n77 Right Arrow\n80 Down Arrow\n```\n\n#### CINT()\n\nThis rounds the contents of the brackets to the nearest integer.\n\n```PRINT CINT(4573.73994596)\n```\n\n4574\n\n#### CIRCLE\n\n```CIRCLE ([X Coordinate], [Y Coordinate]), [Radius], [Colour],[Start],[End],[Aspect]\n```\n\nLets the programmer display a circle. Like all graphics commands, it must be used with the SCREEN command.\n\n#### CLEAR\n\n```CLEAR\n```\n\nResets all variables, strings, arrays and closes all files. The reset command on QBasic.\n\n#### CLOSE\n\n```CLOSE\n```\n\nCloses all open files\n\n```CLOSE #2\n```\n\nCloses only the file opened as data stream 2. Other files remain open\n\n#### CLS\n\n```CLS\n```\n\nClears the active screen. Erases all text, graphics, resets the cursor to the upper left (1,1), and also applies the current background color (this has to be set using the COLOR command) to the whole screen.\n\n#### COLOR\n\n```COLOR [Text Colour], [Background Colour]\n```\n\nThis lets you change the colour of the text and background used when next 'printing' to the current output window. It can be done like this:\n\n```COLOR 14, 01\nPRINT \"Yellow on Blue\"\n```\n\nYou have a choice of sixteen colours:\n\n```00: Black 08: Dark Grey\n01: Dark Blue 09: Light Blue\n02: Dark Green 10: Light Green\n03: Dark Cyan 11: Light Cyan\n04: Dark Red 12: Light Red\n05: Dark Purple 13: Magenta\n06: Orange Brown 14: Yellow\n07: Grey 15: White\n```\n\nThese values are the numbers that you put in the COLOR command.\n\nNote Only screen modes 0, 7, 8, 9, 10 support a background color. To 're-paint' the whole screen in a background colour, use the CLS command.\n\n#### COMMON\n\nDeclares a variable as 'global', which allows its value to be accessed across multiple QBasic programs / scripts (see also the CHAIN command)\n\n```COMMON SHARED [variablename]\n```\n\nEach program that declares 'variablename' as COMMON will share the same value.\n\nNOTE. All COMMON statements must appear at the start of the program (i.e. before any executable statements).\n\n#### CONST\n\nFixes a variable so it can not be changed within the program.\n\n```CONST (name) {AS (type = INTEGER / STRING)} = (expression or value)\n```\n\nFor example :-\n\n```CONST PI = 3.14159265\n```\n\nAssigns the value 3.14159265 to PI.\n\n```CONST PI2 = 2 * PI\n```\n\nPI must be assigned a value before it is used to calculate PI2. Typically all CONST are declared at the beginning of a program.\n\n#### DATA\n\n```DATA [constant]\n```\n\nUse in conjunction with the READ and RESTORE command. Mostly used in programs dealing with graphics, this command lets QBasic read a large number of constants. The READ command accesses the data while the RESTORE command \"refreshes\" the data, allowing it to be used again.\n\n#### DATE\\$\n\nA system variable that always contains the current date as a string in mm-dd-yyyy format. Use it like this:\n\n```a\\$ = DATE\\$\n```\n\n#### DEF SEG\n\n```DEG SEG [=address]\n```\n• address is a segment address that can contain a value of 0 through 65535.\n\nIf address is omitted, DEF SEG resets the current segment address to the default data segment. DEF SEG is used by BLOAD, BSAVE, CALL ABSOLUTE, PEEK, and POKE\n\n#### DEST(Only QB64!)\n\n_DEST sets the current write-to page or image. _DEST image_handle sends the destination image to an image of handle stored in long variable image_handle. _DESt 0 sends the destination image to the current screen being used.\n\n#### DIM\n\nThis is used to declare an array (early versions of QBasic required all variables to be defined, not just arrays greater than 10)\n\n```DIM [Array Name] ([count],[count], ..)\n```\n\nThe Array name can be of any type (Integer, Double, String etc). If not declared, single precision floating point is assumed. Strings can be 'declared' using \\$ sign (Integers with the '%' sign). The QBASIC interpreter tolerates numeric arrays of up to 10 count without these needing to be declared.\n\nNOTE Early versions of QBasic did not explicitly set the contents of an array to zero (see CLEAR command)\n\n```DIM table%(100,2)\n```\n\nCreate an integer array called table% containing 100x2 = 200 entries\n\n```DIM form\\$(5)\n```\n\nCreate a string array called form\\$ containing 5 strings\n\n```DIM quotient(20) AS DOUBLE\n```\n\nCreate an array called quotient that contains 20 double precision numbers\n\n#### DO .. LOOP\n\n```DO\n[program]\nLOOP UNTIL [test condition becomes TRUE]\n```\n\nUsed to create a loop in the program. The [condition] is tested only after the [program] code is executed for the first time (see also WHILE). For example:\n\n```num\\$ = 1\nsum\\$ = 0\nDO\nsum\\$ = 2 * num\\$\nPRINT sum\\$\nnum\\$ = num\\$ + 1\nLOOP UNTIL num\\$ = 13\n```\n\nThis does not work But the following does\n\n```num = 1\nsum = 0\nDO\nsum = 2 * num\nPRINT sum\nnum = num + 1\nLOOP UNTIL num = 13\n```\n\n#### DRAW\n\n```DRAW \"[string expression]\"\n```\n\nUsed to draw a straight line from the current 'cursor' position in the current colour. DRAW defines the direction (up, down etc.) and the length of the line (in pixels). For example:-\n\n```SCREEN 7\nPSET (50, 50), 4\nDRAW \"u50 r50 d50 l50\"\n\n```\n\nThe letter in front of each number is the direction:\n\n```U = Up E = Upper-right\nD = Down F = Lower-right\nL = Left G = Lower-left\nR = Right H = Upper-left\n```\n\nThe drawing 'cursor' is left at the position where the line ends. u50 draws from 50,50 upwards ending at 50,0 r50 draws from 50,0 to the right, ending at 100,0 d50 draws from 100,0 downwards, ending at 100,50 l50 draws from 100,50 to the left, ending at 50,50\n\nThe example shown will thus draw a red 'wire frame' square.\n\nNote: The diagonal from 0,0 to 100,100 will be 100 * root(2) pixels long (i.e. 141)\n\n#### END\n\n```END\n```\n\nSignifies the end of the program. When QBasic sees this command it usually comes up with a statement saying: \"Press Any Key to Continue\".\n\n#### ENVIRON\n\n```ENVIRON [string expression]\n```\n\nNOTE: If you are running QBasic on a Windows system, you will not be able to use this command.\n\nThis command helps you set an environment variable for the duration of the session. On exit from the QBasic.exe interpreter, the variables revert to their original values.\n\n#### EOF()\n\nThis checks if there are still more data values to be read from the file specified in (). EOF() returns a boolean / binary value, a one or zero. 0 if the end of file has not been reached, 1 if the last value in the file has been read (see also LINE INPUT)\n\n```OPEN File.txt FOR INPUT AS #1\nDO\nINPUT #1, text\\$\nPRINT text\\$\nLOOP UNTIL EOF(1)\nEND\n```\n\nNote that, since the INPUT is executed before UNTIL is reached, File.txt must contain at least one line of text - if the file is empty, you will receive an 'ERROR (62) Input past end of file'.\n\n#### ERASE\n\n```ERASE [arrayname] [,]\n```\n\nUsed to erase all dimensioned arrays.\n\n#### ERROR\n\nSystem variable holding a numeric value relating to the processing of the previous line of code. If the line completed without error, ERROR is set to 0. If the line failed, ERROR is set to one of the values shown below. Most commonly used to redirect program flow to error handling code as in :-\n\n```ON ERROR GOTO [line number / label]\n```\n\nIf ERROR is non=zero, program flow jumps to the line number or label specified. If ERROR is zero, program flow continues with the next line below.\n\nTo manually test your program and check to see if the error handling routine runs OK, ERROR can be set manually :-\n\n```ERROR [number]\n```\n\nSet ERROR = number\n\nThe error numbers are as follows:\n\n```1 NEXT without FOR 39 CASE ELSE expected\n2 Syntax Error 40 Variable required\n3 RETURN without GOSUB 50 FIELD overflow\n4 Out of DATA 51 Internal error\n5 Illegal function call 52 Bad file name or number\n7 Out of memory 54 Bad file mode\n8 Label not defined 55 File already open\n9 Subscript out of range 56 FIELD statement active\n10 Duplicate definition 57 Device I/O error\n11 Division by zero 58 File already exists\n12 Illegal in direct mode 59 Bad record length\n13 Type mismatch 61 Disk full\n14 Out of string space 62 Input past end of file\n16 String formula too complex 63 Bad record number\n17 Cannot continue 64 Bad file name\n18 Function not defined 67 Too many files\n19 Yes RESUME 68 Device unavailable\n20 RESUME without error 69 Communication-buffer overflow\n24 Device timeout 70 Permission denied\n25 Device Fault 71 Disk not ready\n26 FOR without NEXT 72 Disk-media error\n27 Out of paper 73 Advanced feature unavailable\n29 WHILE without WEND 74 Rename across disks\n30 WEND without WHILE 75 Path/File access error\n35 Subprogram not defined\n37 Argument-count mismatch\n38 Array not defined\n```\n\nNote that ERROR is set when execution fails, not when the code is 'read' - so, for example, a 'divide by 0' will be found before the result is assigned to a non-existent array variable or written to a non-existent file.\n\n#### EXIT\n\nAllows the immediate exit from a subroutine or a loop, without processing the rest of that subroutine or loop code\n\n```EXIT DEF\n```\n\nExits from a DEF FN function.\n\n```EXIT DO\n```\n\nExits from a DO loop, execution continues with the command directly after the LOOP command\n\n```EXIT FOR\n```\n\nExits from a FOR loop, execution continues with the command directly after the NEXT command\n\n```EXIT FUNCTION\n```\n\nExits a FUNCTION procedure, execution continues with the command directly after the function call\n\n```EXIT SUB\n```\n\nExits a SUB procedure.\n\n#### FOR .. NEXT\n\n```FOR [variable name] = [start value] TO [end value] {STEP n}\n[program code]\nNEXT [variable name]\n```\n\nThe variable is set to the [start value], then program code is executed and at the Next statement the variable is incremented by 1 (or by the STEP value, if any is specified). The resulting value is compared to the [end value] and if not equal program flow returns to the line following the FOR statement.\n\nFor example:\n\n```FOR a = 200 TO 197 STEP-1\nPRINT a\nNEXT a\n```\n\n200 199 198\n\nCare must be taken when using STEP, since it is quite possible to step past the (end value) with the result that the FOR loop will run 'for ever' (i.e. until the user aborts the interpreter or an error occurs), for example :-\n\n```FOR a = 200 TO 197 STEP-2\nPRINT a\nNEXT a\n```\n\n200 198 196 194 192 ... 0 -2 -4 ... -32768 ERROR overflow\n\n#### GOSUB\n\n```GOSUB [subroutine line number / label]\n```\n\nCommand processing jumps to the subroutine specified. When the RETURN command is encountered, processing returns to this point and continues with the line below the GOSUB.\n\n#### IF\n\n```IF [variable or string] [operator] [variable or string] THEN [command] {ELSE [command]}\n```\n\nCompares variables or strings. For example, if you wanted to examine whether or not a user-entered password was the correct password, you might enter:\nWhere a\\$ is the user entered password. Some operators include:\n\"=\"- equal to\n\"<\"- less than (only used when variable or string is a number value)\n\">\"- greater than (only used when variable or string is a number value)\n\"<>\"- does not equal\n\"<=\"- less than or equal to (only used when variable or string is a number value)\n\">=\"- greater than or equal to (only used when variable or string is a number value)\nOne can also preform actions to number values then compare them to other strings or variables using the if command, such as in the below examples:\nIF a+5 = 15 THEN PRINT \"Correct\"\nIF a*6 = b*8 THEN PRINT \"Correct\"\n\n#### INCLUDE (QUICKbasic Only)\n\nQUICKBasic supports the use of include files via the \\$INCLUDE directive:\n\n```(Note that the Qbasic interpreter does NOT support this command.)\n```\n```'\\$INCLUDE: 'foobar.bi'\n```\n\nNote that the include directive is prefixed with an apostrophe, dollar, and that the name of the file for inclusion is enclosed in single quotation mark symbols.\n\n#### INKEY\\$\n\n```[variable] = INKEY\\$\n```\n\nThis is used when you want a program to function with key input from the keyboard. Look at this example on how this works:\n\n```a\\$ = INKEY\\$\nPRINT \"Press Esc to Exit\"\nEND IF a\\$ = CHR\\$(27)\n```\n\nYou can use this in conjunction with the CHR\\$ command or type the letter (e.g. A).\n\n#### INPUT\n\n```INPUT [String Literal] [,or;] [Variable]\n```\n\nDisplays the String Literal, if a semi colon follows the string literal, a question mark is displayed, and the users input until they hit return is entered into the variable. The variable can be a string or numeric. If a user attempts to enter a string for a numeric variable, the program will ask for the input again. The String Literal is option. If the string literal is used, a comma (,) or semicolon (;) is necessary.\n\n#### INPUT #\n\n```INPUT #n [String Literal] [,or;] [Variable]\n```\n\n```INPUT #1, a\\$, b\\$, n, m\n```\n\nReads 4 values from the file that is OPEN as #1. a\\$ is assigned all text until a ',' (comma) or end of line is reached, b\\$ the next segment of text, then two numeric values are interpreted and assigned to n and m.\n\nNote that, within the file, numbers can be separated by 'anything' - so, if a number is not found (for 'n' or 'm') on the current 'line' of the file, the rest of the file will be searched until a number is found. Input is then left 'pointing' at the position in the file after the last number required to satisfy the input statement (see also 'seek #' command)\n\n#### INSTR\n\n```INSTR (start%, Search\\$, Find\\$)\n```\n\nReturns the character position of the start of the first occurrence of Find\\$ within Search\\$, starting at character position 'start%' in Search\\$. If Find\\$ is not found, 0 is returned. start% is optional (default = 1, the first character of Search\\$)\n\n```Pos = INSTR (\"abcdefghi\", \"de\")\n```\n\nreturns 4\n\n#### LEFT\\$()\n\n```A\\$ = LEFT\\$(B\\$,N)\n```\n\nA\\$ is set to the N left most characters of B\\$.\n\n```A\\$ = LEFT\\$(\"Get the start only\",6)\n```\n\nreturns \"Get th\"\n\n#### LET\n\n```LET [variable] = [value]\n```\n\nEarly versions of the QBasic.exe command interpreter required use of the 'LET' command to assign values to variables. Later versions did not.\n\n```LET N = 227 / 99\nLET A\\$=\"a line of simple text\"\n```\n\nis equliavent to :-\n\n```N = 227 / 99\nA\\$=\"a line of simple text\"\n```\n\n#### LINE\n\n```LINE ([X], [Y]) - ([X], [Y]), [Colour Number]\n```\n\nUsed for drawing lines in QBasic. The first X and Y are used as coordinates for the beginning of the line and the second set are used for coordinating where the end of the line is. You must put a SCREEN command at the beginning of the program to make it work.\n\nNote. When in SCREEN 13, the Colour Number == the Palette number\n\n#### LINE INPUT #\n\n```LINE INPUT #1, a\\$\n```\n\nReads a complete line as text characters from the file OPEN as stream #1 and places it in a\\$.\n\nTo find the 'end of line', the QBasic interpreter seaches for the 'Carriage Return' + 'Line Feed' (0x0D, 0x0A) characters. When reading text files created on UNIX/LINUX systems (where the 'Line feed' 0x0A is used on it's own to signify 'end of line'), LINE INPUT will not recognise the 'end of line' and will continue to input until the end of file is reached. For files exceeding 2k characters, the result is an \"Out of String Space\" Error as a\\$ 'overflows'. One solution is to use a text editor able to handle UNIX files to open and 'save as' before attempting to process the file using QBasic.\n\n(NOTE! The commands in this section refer to a third-party program called \"QB64\". Neither QUICKbasic nor Qbasic support _LOADIMAGE, _NEWIMAGE, OR _PUTIMAGE commands. Both QUICKbasic and Qbasic have a \"SCREEN\" command, but it works diffently in those languages than in QB64.)\n\nShows an image. Must be used with the commands SCREEN, _NEWIMAGE and _PUTIMAGE.\n\nExample:\n\nDIM rabbit AS LONG SCREEN _NEWIMAGE(800, 600, 32) rabbit = _LOADIMAGE(\"rabbit.jpg\") _PUTIMAGE (100,100), rabbit\n\n#### LOOP\n\n```DO\n[Program]\nLOOP UNTIL [condition]\n```\n\nUsed to create a loop in the program. This command checks the condition after the loop has started. This is used in conjunction with the DO command.\n\n#### LPRINT\n\n```LPRINT [statement or variable]\n```\n\nPrints out text to a printer. The LPRINT command expects a printer to be connected to the LPT1(PRN) port. If a printer is not connected to LPT1, QBasic displays a \"Device fault\" error message.\n\nIf your printer is connected to a COM port instead, use the MS-DOS MODE command to redirect printing from LPT1 to COMx (for example, to redirect to COM1, use the following command:\n\n``` MODE LPT1=COM1\n```\n\nIf you need to cancel the redirection when finished, use the following command:\n\n``` MODE LPT1\n```\n\n#### MID\\$\n\n``` a\\$=MID\\$(string\\$,start%[,length%])\nMID\\$(string\\$,start%[,length%])=b\\$\n```\n\nIn the first use, a\\$ is set to the substring taken from string\\$ strating with character start% taking Length% characters. If length% is omitted, the rest of the line (i.e. start% and all the characters to the right) are taken.\n\nIn the second use, length% characters of string\\$ are replaced by b\\$ starting at start%. If length% is omitted, the rest of the line is replaced (i.e. start% and all the characters to the right)\n\n#### MOD\n\n``` a MOD b\n```\n\nReturns the remainder of an integer divide of a by b\n\nFor example, 10 MOD 3 returns 1\n\n#### NEWIMAGE(Only QB64!)\n\n_NEWIMAGE is used to set a long variable as the screen dimensions, or can be used with the SCREEN command (See later in Appendix) to directly set the screen details. It is very useful as you can enlarge the SCREEN mode '13' which has RGB color settings if you find the default size too small.\n\nSyntax: _NEWIMAGE(width,length,screen_mode)\n\n• width and length are long variables, while screen_mode is the screen mode format you wish to change.\n\nLike,\n\n``` _NEWIMAGE(1000,1000,24),256\n```\n\nwhere, 256 is the amount of colours\n\nIt is also used to prepare the window screen surface for the image you want to put (first load it using LOADIMAGE).\n\n#### OPEN\n\n```OPEN \"[(path)\\8.3 file name.ext]\" (FOR {INPUT/OUTPUT} AS #{n})\n```\n\nThis opens a file. You have to give the DOS file name, for example:\n\n```OPEN \"data.txt\" FOR INPUT AS #1\n```\n\nOpens the existing file data.txt for reading as data stream #1. Since no path is specified, the file must be in the same folder as the QBasic.exe - if not, processing halts with a 'file not found' error\n\n```OPEN \"C:\\TEMP\\RUN.LOG\" FOR OUTPUT AS #2\n```\n\nOpens an empty file named RUN.LOG in the C:\\TEMP folder for writing data stream #2. Any existing file of the same name is replaced.\n\n#### PALETTE\n\n```PALETTE[palette number, required colour]\n```\n\nFor VGA (SCREEN mode 13) only, sets the Palette entry to a new RGB color. The palette number must be in the range 1-256. The required colour is a LONG integer created from the sum of (required Blue * 65536) + (required Green * 256) + required Red.\n\n#### RANDOMIZE\n\n```RANDOMIZE TIMER\nA = INT((RND * 100)) + 1\n```\n\nRANDOMIZE will set the seed for QBasic's random number generator. With QBasic, it's standard to simply use RANDOMIZE TIMER to ensure that the sequence remains the same for each run.\n\nThe example is a mathematical operation to get a random number from 1 to 100.\n\nINT stands for Integer, RND for Random and \"*\" stands for the limit upto which the random number is to be chosen. The \"+ 1\" is just there to ensure that the number chose is from 1 to 100 and not 0 to 99.\n\nNote: Subsequent calls of this function do not guarantee the same sequence of random numbers.\n\n```READ AIM(I)\n```\n\nUsed in conjunction with the DATA command, this command lets QBasic read data. This is mostly used when dealing with large quantities of data like bitmaps.\n\n#### REM or '\n\n```REM {comments}\n```\n\nWhen the interpreter encounters REM or \" ' \" (a single quote) at the start of a line, the rest of the line is ignored\n\n#### RETURN\n\n```RETURN\n```\n\nSignifies that it is the end of a subroutines\n\n#### RND\n\n```RANDOMIZE TIMER\nA = INT((RND * 100)) + 1\n```\n\nRND will provide a random number between 0 and 1.\n\nThe example is a mathematical operation to get a random number from 1 to 100. RANDOMIZE TIMER will set the initial seed to a unique sequence. INT stands for Integer, RND for Random and \"*\" stands for the limit upto which the random number is to be chosen. The \"+ 1\" is just there to ensure that the number chose is from 1 to 100 and not 0 to 99.\n\nInternally, the seed a 24-bit number, iterated in the following method: rnd_seed = (rnd_seed*16598013+12820163) MOD 2^24\n\n#### PLAY\n\n```PLAY \"[string expression]\"\n```\n\nUsed to play notes and a score in QBasic on the PC speaker. The tones are indicated by letters A through G. Accidentals are indicated with a \"+\" or \"#\" (for sharp) or \"-\" (for flat) immediately after the note letter. See this example:\n\n```PLAY \"C C# C C#\"\n```\n\nWhitespaces are ignored inside the string expression. There are also codes that set the duration, octave and tempo. They are all case-insensitive. PLAY executes the commands or notes the order in which they appear in the string. Any indicators that change the properties are effective for the notes following that indicator.\n\n```Ln Sets the duration (length) of the notes. The variable n does not indicate an actual duration\namount but rather a note type; L1 - whole note, L2 - half note, L4 - quarter note, etc.\n(L8, L16, L32, L64, ...). By default, n = 4.\nFor triplets and quintets, use L3, L6, L12, ... and L5, L10, L20, ... series respectively.\nThe shorthand notation of length is also provided for a note. For example, \"L4 CDE L8 FG L4 AB\"\ncan be shortened to \"L4 CDE F8G8 AB\". F and G play as eighth notes while others play as quarter notes.\nOn Sets the current octave. Valid values for n are 0 through 6. An octave begins with C and ends with B.\nRemember that C- is equivalent to B.\n< > Changes the current octave respectively down or up one level.\nNn Plays a specified note in the seven-octave range. Valid values are from 0 to 84. (0 is a pause.)\nCannot use with sharp and flat. Cannot use with the shorthand notation neither.\nMN Stand for Music Normal. Note duration is 7/8ths of the length indicated by Ln. It is the default mode.\nML Stand for Music Legato. Note duration is full length of that indicated by Ln.\nMS Stand for Music Staccato. Note duration is 3/4ths of the length indicated by Ln.\nPn Causes a silence (pause) for the length of note indicated (same as Ln).\nTn Sets the number of \"L4\"s per minute (tempo). Valid values are from 32 to 255. The default value is T120.\n. When placed after a note, it causes the duration of the note to be 3/2 of the set duration.\nThis is how to get \"dotted\" notes. \"L4 C#.\" would play C sharp as a dotted quarter note.\nIt can be used for a pause as well.\nMB MF Stand for Music Background and Music Foreground. MB places a maximum of 32 notes in the music buffer\nand plays them while executing other statements. Works very well for games.\nMF switches the PLAY mode back to normal. Default is MF.\n```\n\n#### PRINT\n\n```PRINT [Argument] [,or;] [Argument]...\n```\n\nDisplays text to the screen. The Argument can be a string literal, a string variable, a numeric literal or a numeric variable. All arguments are optional.\n\n```PRINT #[n] [,or;] [Argument]...\n```\n\nSaves data to the file that is 'OPEN FOR OUTPUT AS #[n]' or we can use ? symbol for print command\n\n#### PSET\n\n```PSET ([X coordinate],[Y coordinate]), [Pixel Colour]\n```\n\nThis command displays pixels, either one at a time or a group of them at once. For the command to work, the program must have a SCREEN command in it.\n\n#### SCREEN\n\n```SCREEN [Screen Mode Number]\n```\n\nThis command is used for displaying graphics on the screen. There are ten main types of screen modes that can be used in QBasic depending on the resolution that you want. Here is a list of what screen modes you can choose from:\n\nSCREEN 0: Textmode, cannot be used for graphics. This the screen mode that text based programs run on.\n\nSCREEN 1: 320 x 200 Resolution. Four Colours\n\nSCREEN 2: 640 x 200 Resolution. Two Colours (Black and White)\n\nSCREEN 7: 320 x 200 Resolution. Sixteen Colours\n\nSCREEN 8: 640 x 200 Resolution. Sixteen Colours\n\nSCREEN 9: 640 x 350 Resolution. Sixteen Colours\n\nSCREEN 10: 640 x 350 Resolution. Two Colours (Black and White)\n\nSCREEN 11: 640 x 480 Resolution. Two Colours\n\nSCREEN 12: 640 x 480 Resolution. Sixteen Colours\n\nSCREEN 13: 320 x 200 Resolution. 256 Colours. (Recommended)\n\nNote. In SCREEN 13 you have a colour Palette of 256 colours. The PALETTE is pre-set by Windows however you can change the RGB values using the PALETTE command.\n\n#### SEEK\n\n```SEEK #[file number], 1\n```\n\nRepositions the 'input #' pointer to the beginning of the file.\n\n#### SGN\n\n```SGN(expression yielding a numeric value)\n```\n\nYields the 'sign' of a value, -1 if < 0, 0 if 0, 1 if > 0\n\n#### SHELL\n\nThe 'SHELL' command is used in Qbasic to issue a command to Command Prompt/Windows Shell . The 'Shell' command is used along with a string that contains commands that would be understood by any of the above software. The string enclosed commands are much like that of MS-DOS\n\nExample: SHELL can be used with a 'DIR' command to make a directory of files in a certain folder or path.\n\n#### SLEEP\n\n```SLEEP [n]\n```\n\nExecution is suspended for n seconds\n\n#### SOUND\n\n```SOUND [frequency], [duration]\n```\n\nUnlike the BEEP command, this produces a sound from the PC speakers that is of a variable frequency and duration. The frequency is measured in Hertz and has a range from 37 to 32767. Put in one of these numbers in the frequency section. The duration is clock ticks that is defaulted at 18.2 ticks per second.\n\n#### STR\\$\n\nConverts a numeric value into a text (string) character\n\n```A\\$ = STR\\$(expression yielding a numeric value)\n```\n\nThe numeric value is converted into text characters and placed into A\\$. Use to convert numbers into a text string.\n\nWARNINGS.\n\n1) If the result is positive, a leading 'space' is added (STR\\$(123) = \" 123\" and not \"123\" as might be expected). If the result is negative, instead of a space you get a '-' (minus sign), i.e. STR\\$(-123) = \"-123\" and not \" -123\" as might be expected from the positive behaviour.\n\n2) When converting a float (mumb!, numb#) less than 0.1, the string value may be rendered in 'scientific notation', with 'D' used rather than '*10^' (for example \"5.nnnnnnD-02\" rather than \" .05nnnnnn\" or \"5.nnnnnn*10^-02\"). This only occurs when the number of significant digits needs to be preserved (so .03000000 is rendered as \" .03\", whilst .030000001 becomes \" 3.0000001D-02\"), again perhaps not what you might expect.\n\nSee also CHR\\$ for converting an ascii value into a string character.\n\nSee also LEFT\\$, MID\\$, RIGHT\\$ for extracting sub-strings from a line of text.\n\n#### SYSTEM\n\n```SYSTEM\n```\n\nThe .bas exits, the QBasic.exe interpreter is closed and 'control' passes to the Command Window c:\\ prompt (or next line of a calling .cmd script etc.)\n\nNOTE!: This only works when you start your program at the command prompt using the \"/run\" parameter! (EX: \"Qbasic /run MyProg.bas\") Otherwise, Qbasic assumes you opened your program to make changes, and thus \"SYSTEM\" drops you back at the editor screen.\n\n#### THEN\n\n```[Command] [variable] = [value] THEN GOTO [line command value]\n```\n\nUsed in conjunction with the GOTO or IF condition commands. It tells the computer what to do if a certain condition has been met.\n\n#### TO\n\n```[Command] [Variable] = [Value] TO [Value]\n```\n\nUsually used to input a number of variables.\n\n```FOR a = 400 TO 500\nPRINT a\nNEXT a\n```\n\nThis example will print all numbers from 400 to 500. Instead of declaring all values separately, we can get them all declared in one go.\n\n#### USING\n\n```USING \"format\";\n```\n\nUsed to format the output of data from PRINT commands. Normally, the QBasic interpreter will print a number as 8 characters with as many leading spaces as necessary. To change this behavour, the USING command can be used to format the output. For example ..\n\n```IF n > 99 THEN PRINT #1, USING \"###\"; n; ELSE IF n > 9 AND n<=99\n```\n\nTHEN PRINT #1, USING \"0##\"; n; ELSE PRINT #1, USING \"00#\"; n;\n\n.. will output n from 0 to 999 with leading zeros. Note the ';' after the n. This means 'don't start a new line' and results in the next PRINT #1 adding data directly after the comma (',') Qbasic automatically inserts instead of a line.\n\n#### VAL()\n\n```name=VAL([variable\\$])\n```\n\nConverts the [variable string] contents into a numeric value so it can be used in calculations. If (name) is an INTEGER type, the VAL is rounded down. See also STR\\$.\n\nA\\$ = \"2\"\n\nB\\$ = \"3\"\n\nX = VAL(A\\$) + VAL(B\\$)\n\nPRINT A\\$; \" + \"; B\\$; \" =\"; X\n\n#### WHILE ... WEND\n\n```WHILE {NOT} [test condition is true]\n[program code to execute]\nWEND\n```\n\nThe condition is tested and if true (or NOT true) the [program] code is executed until WEND is reached, at which point control passes back to the WHILE line.\n\n```WHILE NOT (EOF(1))\nLINE INPUT #1, A\\$\nPRINT #2, A\\$\nWEND\n```\n\nWhile the end of file #1 has not been reached, read each complete line and write it to file #2.\n\nUnlike FOR and DO, it is not possible to EXIT from a WHILE loop" ]
[ null ]
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https://www.codingninjas.com/codestudio/problems/count-subsequences_630524
[ "Problem title\nDifficulty\nAvg time to solve\n\nCount Ways\nEasy\n24 mins\nTerms Of AP\nEasy\n10 mins\nCount Distinct Element in Every K Size Window\nEasy\n15 mins\nFire in the cells.\nHard\n15 mins\nMerge Two BSTs\nModerate\n10 mins\nSearch For Integers With Given Difference And At Given Distance\nModerate\n20 mins\nWord Break II\nHard\n15 mins\nRecycling Pens\nEasy\n15 mins\nMax Path Value\nModerate\n15 mins\nMaximal AND Subsequences\nModerate\n15 mins", null, "19\n\n# Count Subsequences\n\nDifficulty: MEDIUM\nAvg. time to solve\n30 min\nSuccess Rate\n60%\n\nProblem Statement\n\n#### A subsequence of a given array is an array generated by deleting some elements of the given array with the order of elements in the subsequence remaining the same as the order of elements in the array.\n\n##### Note :\n``````As this value might be large, print it modulo 10^9 + 7\n``````\n##### Input format :\n``````The first line contains an integer 'T' which denotes the number of test cases. Then the test cases follow :\n\nThe first line of each test case contains an integer 'N' representing the size of the array/list.\n\nThe second line contains 'N' single space - separated integers representing the elements of the array.\n``````\n##### Output Format :\n``````For each test case, print the total number of subsequences in which all elements are equal.\n\nThe output of each test case will be printed in a separate line.\n``````\n##### Note :\n``````You don't need to print anything. It has already been taken care of. Just implement the given function.\n``````\n##### Constraints :\n``````1 <= T <= 100\n1 <= N <= 10^5\n0 <= ARR[i] <= 10^9\n\nTime Limit : 1 sec\n``````\n##### Sample Input 1 :\n``````2\n1\n5\n2\n1 0\n``````\n##### Sample Output 1 :\n``````1\n2\n``````\n##### Explanation For Sample Input 1 :\n``````For the first query with input {5}, the total number of subsequences would be 2 which are {{\"\"}, {5}}.\nOut of these subsequences, there is only one subsequence which has all the same elements and that is {5} itself.\n\nFor the second query with input {1, 0}, the total number of subsequences would be 4 which are {{\"\"}, {1}, {0}, {1, 0}}.\nOut of these subsequences, there are two subsequences which have all the same elements and they are {1}, {0}.\n``````\n##### Sample Input 2 :\n``````2\n2\n1 1\n3\n1 1 1\n``````\n##### Sample Output 2 :\n``````3\n7\n``````", null, "", null, "", null, "Console" ]
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https://or.stackexchange.com/questions/2728/how-to-linearize-a-constraint-with-max
[ "# How to linearize a constraint with max\n\nI would like to linearize a constraint with max. I have the following constraint:\n\n$$\\max_{pcj}X_{pwcj}\\leqslant L_{wk}.$$\n\nWith this constraint, I would like to ensure that for $$\\forall w \\in W$$, no matter the values of $$p$$, $$c$$ and $$j$$, the max value of the decision variable $$X_{pwcj}$$ must be less than or equal to the value of the decision variable $$L_{wk}$$.\n\n1. How can I linearize this function? I've seen some examples, such as How to linearize a constraint with a maximum of binary variables times some coefficient in the right-hand-side. However, their decision variable has only one index. In my case, I have 4 indexes and $$w$$ must be fixed at each time.\n\nActually, I would like to check for each set of $$w$$ occurrences in a sequence given by $$j$$:\n\n$$H_w=\\max \\{X_{pwcj}, X_{pwcj+1},\\ldots, X_{pwcn}\\}$$\n\nI was wondering what the best way to write my constraint could be:\n\n$$H_{w}\\leqslant L_{wk}$$ where $$n$$ is the max number of $$j$$ positions, but I am not sure.\n\n1. If I use this notation, I think that I could do something like this: How to formulate (linearize) a maximum function in a constraint? proposed by @LarrySnyder610, however, in my case I have a set of the same variables thus I am not sure that I can do this.\n\nIf this notation can be used in my case, could I ignore the big M value, since both of my variables are binary?\n\nCould someone help me with these questions?\n\nWelcome to OR.SE! If you're looking to enforce $$\\max\\limits_{pcj}X_{pwcj} \\leq L_{wk}, \\ \\forall w,k$$ then simply using the constraint $$X_{pwcj} \\leq L_{wk}, \\ \\forall p,w,c,j,k$$ will do the trick. This will ensure every value (including the maximum) will be at most $$L_{wk}$$.\n• @campioni, thank you for the additional detail. Is the purpose of the constraint to ensure that $X$ is forced to 0 if $L =0$ for that machine (i.e., can't use machine if not located anywhere)? If not, what is the relationship you are trying to establish between $X$ and $L$? Oct 3, 2019 at 14:18" ]
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https://money.stackexchange.com/questions/119822/how-can-you-calculate-the-sweet-spot-for-a-risk-reward-ratio-in-a-partnership
[ "# how can you calculate the sweet spot for a risk / reward ratio in a partnership\n\nLet's take 2 parties, A and B.\n\nA has a tool that can raise the value of an investment by a known quantity (Q) with an estimated risk factor (R). A failure would incur a loss of (L).\n\nB has money (M) to put in the tool.\n\nlet's plug some numbers to illustrate:\n\n• R = 0.8 (80% chance of success)\n• Q = 0.2 (20% profit if successful)\n• L = 0.1 (10% loss if failure)\n• M = \\$100 invested by B\n\nThe following scenarios are possible:\n\n• The process is successful, M becomes M * (1+Q) = \\$120\n• The process is a failure, M becomes M * (1-L) = \\$90\n\nA is not taking any risks, but is providing an essential part of the process. A cannot lose money. B can lose the money put in the process\n\nB can earn \\$20, or loose \\$10. If B earns money, he has to share with A.\n\nHow one could express the ratio at which the money is split using R, Q and L?\n\n• `A is not taking any risks, ... A cannot lose money.` Is A getting income from the partnership while working on the project? Jan 31, 2020 at 23:19\n• In this specific context A is owning a trading software (with known financial value) and B is putting money to use it. Jan 31, 2020 at 23:21\n• A still needs to eat and pay the rent. How's he doing that during the partnership? Feb 1, 2020 at 2:28\n• @RonJohn: A has developed the software at his own cost and is not bound to 'B', he's free to exploit it with other people. In that context we can consider that 'A' made an investment as well and the the value of what he brings to the table grows as it becomes more profitable. Feb 1, 2020 at 10:55\n• I'm not clear exactly what you are asking for. Normally, If someone came to a deal like this one would write a special allocation rule for dividing profits and losses between A and B, that should share some of the profit, if there was any, with A, but not all of it. It would probably be less than 50-50 since B is bearing loss risk that A is not. But lots of factors not in the question could impact the kind of deal that makes sense. Feb 26, 2021 at 6:28\n\nThere's not much to do here:\n\nThink in terms of ten separate plays of \\$100 each and the gain is (8 * \\$20) while the loss is (2 * \\$10) such that the average profit is (\\$140 / 10) and as 14% .\n\nThe percentage that Person A takes is just competitive marketing. Hedge funds take 20% of profit plus 2% of each year's beginning balance. Asset managers take 1% of each year's beginning balance and 0% of profit.\n\nIf Person A could guarantee a 7% profit then Person A could possibly earn 7%. But that's an insurance company endeavor.\n\nThis can be split into two issues: how much is this worth to B, and of that value that A is supplying to B, how much can A claim? The second question is a matter of negotiation, and is going to depend on issues such as supply and demand, how much confidence people have that those probabilities are correct, and other factors such as possibly their relationship.\n\nAs for how much value A is supplying to B, that is difficult to quantify, since you need some risk discounting factor. One method of modeling risk is to assume that the more money a person has, the less each additional dollar adds value. This then justifies using some diminishing returns function to measure the value of a person's net worth. What function to use is somewhat arbitrary, but the log function is often used. Using the log function, we would have that this investment results in some value X, and the log of X is equal to the weighted logs of the possibilities:\n\nlog(X) = 0.8(log(120))+0.2(log(90))\n\nThis gives us that X = 113.29. So this calculation suggests that B is getting \\$13.29 of value from this arrangement. However, it's not quite as simple as saying \"Okay, then A can't charge more than \\$13.29\"; now that we've introduced log, the math isn't linear any more. A could actually charge \\$17 and B would still be getting value from it.\n\nA further complication is that B's investment of \\$100 probably isn't their entire net worth. If they have \\$1000 total, then rather than log(X) = 0.8(log(120))+0.2(log(90)), you'd have log(X) = 0.8(log(1020))+0.2(log(990)), which increases X significantly. But they're probably putting that other money in further investments, so to calculate their total risk, you'd have to look at what those investments are and how much correlation they have with this one.\n\nYou'd also have to look at further factors such as the time value of money." ]
[ null ]
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https://plazaiberica.pl/an/29174_an-aggregate-supply-curve-shows-the.html
[ "# an aggregate supply curve shows the\n\n•", null, "#### Chapter 08 - Questions - SV (1).pdf - Chapter\n\nChapter 8 — Aggregate Demand and Aggregate Supply 8-01. An aggregate demand (AD) curve shows the A. amount of a particular good people are willing and able to buy at a particular price, ceteris paribus. B. real output (Real GDP) people are willing and able to sell at different price levels, ceteris paribus. C. real output (Real (GDP) people are willing and able to buy and to sell at\n\nGet Price\n•", null, "#### Aggregate supply - Economics Help\n\nThe aggregate supply curve shows the amount of goods that can be produced at different price levels. When the economy reaches its level of full capacity (full employment – when the economy is on the production possibility frontier) the aggregate supply curve becomes inelastic because, even at higher prices, firms cannot produce more in the short term ; The aggregate supply curve is related\n\nGet Price\n•", null, "#### Short-Run Aggregate Supply: Meaning, Its\n\n23.04.2021· Short-run aggregate supply. In a graph where the X-axis represents aggregate output, and the Y-axis represents the price level, the short-run aggregate supply (SRAS) curve has an upward slope. It shows an increase in the price level encourages an increase in aggregate output, represented by real GDP. Remember, in the short run, we are assuming\n\nGet Price\n•", null, "#### Macroeconomics Chapter 7 Quiz Flashcards |\n\nAn aggregate supply curve shows the: Answers: A. price level at which real domestic output will be in equilibrium. B. price level at which real domestic output will be purchased. C. level of real domestic output that will be produced at each possible price level. D. level of real domestic output that will be purchased at each possible price level. A. The slope of the immediate-short-run\n\nGet Price\n•", null, "#### Aggregate Supply: Problems 1 | SparkNotes\n\nThe aggregate supply curve shows the relationship between the price level and the quantity of goods and services supplied in an economy. Problem : What is the equation for the aggregate supply curve in the short run? The equation for the aggregate supply curve in the short run is Y = Ynatural + a(P - Pexpected). Problem : What does each of the terms mean in the equation for the short-term\n\nGet Price\n•", null, "#### What is aggregate supply curve? –\n\n26.11.2019· The aggregate supply curve shows the total quantity of output—real GDP—that firms will produce and sell at each price level. Why aggregate supply curve is upward sloping? The short-run aggregate supply curve is upward sloping because the quantity supplied increases when the price rises. As a result, there is a positive correlation between the price level and output, which is shown on the\n\nGet Price\n•", null, "#### 1. The short-run aggregate supply curve shows:\n\n1. The short-run aggregate supply curve shows: a. What happens to output in an economy as the price level changes, holding all other determinants of real GDP constant.\n\nGet Price\n•", null, "#### What causes increases or decreases in aggregate\n\n21.02.2020· The aggregate supply curve shows the total quantity of output—real GDP—that firms will produce and sell at each price level. The graph shows an upward sloping aggregate supply curve. How does inflation affect aggregate supply? When the aggregate supply of goods and services decreases because of an increase in production costs, it results in cost-push inflation. In order to compensate,\n\nGet Price\n•", null, "#### Answered: The short-run aggregate supply\n\nThe short-run aggregate supply curve shows: O The relationship between the price level and aggregate expenditure O What happens to the level of real GDP suppliers are willing and able to produce in an economy as the overall price level changes, during a period in which output prices can change but input prices are fixed How firms respond to changes in interest rates O What happens to\n\nGet Price\n•", null, "#### Aggregate demand and aggregate supply -\n\n17.11.2014· The aggregate supply curve shows the total quantity of output—real GDP—that firms will produce and sell at each price level. The graph below shows an aggregate supply curve. Let's begin by walking through the elements of the diagram one at a time: the horizontal and vertical axes, the aggregate supply curve itself, and the meaning of the potential GDP vertical line. The aggregate supply\n\nGet Price\n•", null, "#### Unit 6 - 1 Aggregate Supply curve shows the\n\nAggregate Supply curve shows the relationship between the price level and the real GDP supplied in an economy. a. Under what circumstances will the AS curve have a flat segment? The AS Curve will have a flat segment when the price of factors of production is fixed, with little or no upward pressure on price. b. When an economy has a vertical AS curve? An economy will have a vertical AS curve\n\nGet Price\n•", null, "#### Solved > 11) The aggregate supply curve\n\n11) The aggregate supply curve shows the relationship between. A) potential GDP and the price level. B) potential GDP and real GDP. C) the quantity of real GDP supplied and the price level. D) the quantity of real GDP supplied and the interest rate. E) potential GDP and the aggregate demand curve.\n\nGet Price\n•", null, "#### Answered: 1. The short-run aggregate supply\n\n1. The short-run aggregate supply curve shows: a. The relationship between the price level and aggregate expenditure b. What happens to the level of real GDP suppliers are willing and able to produce in an economy as the overall price level changes, during a period in which output prices can change but input prices are fixed c.\n\nGet Price\n•", null, "#### Interpreting the AD-AS Model | Macroeconomics\n\nEquilibrium in the Aggregate Demand–Aggregate Supply Model. Figure 1 combines the AS curve and the AD curve from Figures 1 & 2 on the previous page and places them both on a single diagram. The intersection of the aggregate supply and aggregate demand curves shows the equilibrium level of real GDP and the equilibrium price level in the\n\nGet Price\n•", null, "#### The short run aggregate supply curve shows\n\nThe short-run aggregate supply curve shows the various amounts of real output that producers are willing to a. sell at different profit levels. b. sell at different price levels. c. buy at different income levels. d. buy at different price levels. Refer to Exhibit 8-3.\n\nGet Price\n•", null, "#### Solved: An Aggregate Supply Curve Shows The:\n\nAn aggregate supply curve shows the: level of real domestic output that will be produced at each possible price level. level of real domestic output that will be purchased at each possible price level. price level at which real domestic output will be purchased price\n\nGet Price\n•", null, "#### the aggregate supply curve - am-lift.de\n\nThe aggregate supply curve shows the total supply in an economy at different price levels. Generally, the aggregate supply curve slopes upwards a higher price level encourages firms to supply more. However, there are different possible slopes for the aggregate supply curve. It Read More. Aggregate supply Wikipedia. In the standard aggregate supplyaggregate demand model, real output\n\nGet Price\n•", null, "#### BUSI2003-S.Bagheri_UEx# 6.docx - BUSI2003\n\nAggregate Supply curve shows the relationship between the price level and the real GDP supplied in an economy. a. Under what circumstances will the AS curve have a flat segment? The AS curve will have a flat segment when the price factors of production is fixed, with little or no upward rise in price for production costs, while level of production increases. This happens because of the fixed\n\nGet Price\n•", null, "#### Chapter 08 - Questions - SV (1).pdf - Chapter\n\nChapter 8 — Aggregate Demand and Aggregate Supply 8-01. An aggregate demand (AD) curve shows the A. amount of a particular good people are willing and able to buy at a particular price, ceteris paribus. B. real output (Real GDP) people are willing and able to sell at different price levels, ceteris paribus. C. real output (Real (GDP) people are willing and able to buy and to sell at\n\nGet Price\n•", null, "#### What causes the aggregate supply curve to shift\n\n18.06.2020· The aggregate supply curve shows the total quantity of output—real GDP—that firms will produce and sell at each price level. The graph shows an upward sloping aggregate supply curve. 27 Related Question Answers Found What shifts the aggregate supply curve? Reasons for Shifts The short-run aggregate supply curve is affected by production costs including taxes, subsidies, price of\n\nGet Price", null, "" ]
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https://socoder.net/?Showcase=40790&DockBar=1
[ "-=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- (c) WidthPadding Industries 1987 0|584|0 -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=- -=+=-\nSoCoder -> Showcase Home -> Arcade\n\ntherevillsgames", null, "Created : 29 June 2014\nEdited : 29 June 2014\nSystem : Android\nLanguage : Monkey\n\n### Don't Tap The Zombies\n\nAndroid\nHTML5\nScreenshots", null, "A quick little game about not touching the Zombies!\n\nDon't Tap The Zombies is an addictive survival game.\nWatch your step and don't tap any Zombies!\n\nTwo Game Modes:\n* Sprint 25 / Sprint 50\n* Survival Mode / Fast\n\nIn Sprint mode you need to tap 25/50 safe tiles as fast as you can in order.\nIn Survival mode you need to tap safe tiles for as long as you can in order.\n\nYou can collect guns, which will gives you the ability to shoot Zombies, but you must still tap the safe tile!\n\nAndroid Version:\n\nHTML5 Version:\nhttps://www.monkey-x.com/Community/posts.php?topic=8688&app_id=319", null, "", null, "", null, "" ]
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https://www.geeksforgeeks.org/maximum-array-least-twice-elements/?ref=rp
[ "Maximum in array which is at-least twice of other elements\n\n• Difficulty Level : Basic\n• Last Updated : 09 Apr, 2021\n\nGiven an array of integers of length n. Our task is to return the index of the max element if the it is at least twice as much as every other number in the array. If the max element does not satisfy the condition return -1.\nExamples:\n\nInput : arr = {3, 6, 1, 0}\nOutput : 1\nHere, 6 is the largest integer, and for\nevery other number in the array x, 6 is\nmore than twice as big as x. The index of\nvalue 6 is 1, so we return 1.\n\nInput : arr = {1, 2, 3, 4}\nOutput : -1\n4 isn't at least as big as twice the value\nof 3, so we return -1.\n\nApproach : Scan through the array to find the unique largest element m, keeping track of it’s index maxIndex. Scan through the array again. If we find some x != m with m < 2*x, we should return -1. Otherwise, we should return maxIndex.\n\nC++\n\n // CPP program for Maximum of// the array which is at least// twice of other elements of// the array.#includeusing namespace std;     // Function to find the// index of Max element// that satisfies the// conditionint findIndex(int arr[], int len) {         // Finding index of    // max of the array    int maxIndex = 0;    for (int i = 0; i < len; ++i)        if (arr[i] > arr[maxIndex])            maxIndex = i;         // Returns -1 if the    // max element is not    // twice of the i-th    // element.    for (int i = 0; i < len; ++i)                    if (maxIndex != i &&            arr[maxIndex] < 2 * arr[i])            return -1;         return maxIndex;}     // Driver functionint main(){    int arr[] = {3, 6, 1, 0};    int len = sizeof(arr) / sizeof(arr);         cout<<(findIndex(arr, len));} // This code is contributed by Smitha Dinesh Semwal\n\nJava\n\n // Java program for Maximum of the array// which is at least twice of other elements// of the array.import java.util.*;import java.lang.*; class GfG {         // Function to find the index of Max element    // that satisfies the condition    public static int findIndex(int[] arr) {                        // Finding index of max of the array        int maxIndex = 0;        for (int i = 0; i < arr.length; ++i)            if (arr[i] > arr[maxIndex])                maxIndex = i;                 // Returns -1 if the max element is not        // twice of the i-th element.               for (int i = 0; i < arr.length; ++i)                        if (maxIndex != i && arr[maxIndex] < 2 * arr[i])                return -1;                 return maxIndex;    }         // Driver function    public static void main(String argc[]){        int[] arr = new int[]{3, 6, 1, 0};        System.out.println(findIndex(arr));    }}\n\nPython3\n\n # Python 3 program for Maximum of# the array which is at least twice # of other elements of the array. # Function to find the index of Max# element that satisfies the conditiondef findIndex(arr):             # Finding index of max of the array    maxIndex = 0    for i in range(0,len(arr)):        if (arr[i] > arr[maxIndex]):            maxIndex = i         # Returns -1 if the max element is not    # twice of the i-th element.        for i in range(0,len(arr)):            if (maxIndex != i and                arr[maxIndex] < (2 * arr[i])):            return -1             return maxIndex          # Driver codearr = [3, 6, 1, 0]print(findIndex(arr)) # This code is contributed by Smitha Dinesh Semwal\n\nC#\n\n // C# program for Maximum of the array// which is at least twice of other elements// of the array.using System; class GfG {         // Function to find the index of Max element    // that satisfies the condition    public static int findIndex(int[] arr) {                     // Finding index of max of the array        int maxIndex = 0;        for (int i = 0; i < arr.Length; ++i)            if (arr[i] > arr[maxIndex])                maxIndex = i;                 // Returns -1 if the max element is not        // twice of the i-th element.        for (int i = 0; i < arr.Length; ++i)                    if (maxIndex != i && arr[maxIndex] < 2 * arr[i])                return -1;                 return maxIndex;    }         // Driver function    public static void Main()    {        int[] arr = new int[]{3, 6, 1, 0};        Console.WriteLine(findIndex(arr));    }} // This code is contributed by vt_m.\n\nPHP\n\n \\$arr[\\$maxIndex])            \\$maxIndex = \\$i;         // Returns -1 if the    // max element is not    // twice of the i-th    // element.    for (\\$i = 0; \\$i < \\$len; ++\\$i)                    if (\\$maxIndex != \\$i and            \\$arr[\\$maxIndex] < 2 * \\$arr[\\$i])            return -1;         return \\$maxIndex;}     // Driver Code\\$arr = array(3, 6, 1, 0);\\$len = count(\\$arr); echo findIndex(\\$arr, \\$len); // This code is contributed by anuj_67.?>\n\nJavascript\n\n\n\nOutput:\n\n1\n\nTime Complexity", null, "My Personal Notes arrow_drop_up" ]
[ null, "https://www.geeksforgeeks.org/wp-content/ql-cache/quicklatex.com-e8a7d26c0a94af32b616e4bcc90839db_l3.png", null ]
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https://forum.arduino.cc/index.php?topic=142609.0
[ "Go Down\n\nTopic: BPM (Beats per Minute) detection program - my version (Read 19834 times)previous topic - next topic\n\ndreggory", null, "Jan 14, 2013, 05:49 amLast Edit: May 19, 2015, 10:06 pm by dreggory\nOkay so, after many months of study and tinkering, writing and re-writing code I finally made a working(at least with version 22)(see my last post I uploaded the library and my code that works with IDE 1.6)  beat detection program it takes 10 seconds to get a rough beat and then starts honing in after 20 more seconds if I remember correctly... it's been a while since I messed with it. I found a much cooler solution to my problem and no longer need this code, but I put the absolute peak of my focus and effort into it for many months so it seems a waste to not share it.\n\nbe nice to me I'm not a programmer and I only had one class of java in college. but I know enough to know that it's sloppy, not very well documented and downright confusing at some parts. for example, I start using some weird variables like \"fruit\" or \"Fruity\" my brain was fried and I needed a variable with a unique name... sorry. Despite all this I'll do my best to answer any questions.\n\nkeep in mind that beat detection is hard enough to do for a person and their instincts let alone to define it in an algorithm but this is still fun to see it respond to faster or slower beats even though it might be off a little (it's not as good as we all hope it was.) sometimes its dead on though.\n\nwhen I started this project I knew nothing of Arduinos or micro controllers for that matter, but after lurking 100% (out of fear of angering grumpy Arduino community members) I learned a lot. in fact I feel that everyone needs to try and figure things out without making things a problem for someone else or at least until you are really stuck.\n\nI still don't feel like I belong posting anything on here cause I don't know very much. but I wanted to contribute this so here it is\n\noh and this is actually a two Arduino project, one for detecting the BPM and one for controlling animations on an 8X8 RGB LED matrix. if anyone is interested in the animation side of this project I'll post it in the appropriate section and add a link to it in this post.\nI fully intended on cleaning up this code and making it pretty and easy to understand but I am busy with other projects, a three year old, job, new business and life in general so this is what you get, sorry.\n\nsome things you'll need besides my code:\n\nSimpleTimer library\n\nyou will also need this web site as reference as to how to hook up the MSGEQ7 chip for audio data input:\n\nhttp://nuewire.com/info-archive/msgeq7-by-j-skoba/\n\ndreggory", null, "#1\nJan 14, 2013, 05:54 amLast Edit: Jan 14, 2013, 06:22 am by dreggory Reason: 1\nand finally my sketch:\nCode: [Select]\n/* Beat Detection Program\n\nthis program uses a MSGEQ7 chip to collect sound data,\nthe chip seperates the audio into 7 bands centered on 63hz, 160hz, 400hz, 1khz,\n2.5khz, 6.25khz and 16khz(which from my tests is more like 15khz)\n\nthen it takes the average of however many samples are in 1.25 seconds\nand compares it to the current sample. If the current sample is higher\nthan the average multiplied by C, it counts a beat.\n\nit takes the average amounts of beats in ten seconds and converts that\nto beats per minute.\n\nIt then takes 60 divided by the average beats per minute and multiplies\nit by 1000 to convert beats per minute into a delay to use between\nanimations.\n\nThe animations are done on an eight by eight LED Matrix. I'm using four of the\n74hc595 ic shift registers to extend the outputs of the Arduino Uno.\nI used common cathode RGB LEDs so if the binary animations seem funny\nit may be that you have common anode LEDs.\n\n*/\n\n#include <Wire.h>\n\n#include <SimpleTimer.h>\nSimpleTimer timer;\n\n// these are for the MSGEQ7\nint analogPin = 0; // read from multiplexer using analog input 0\nint strobePin = 2; // strobe is attached to digital pin 2\nint resetPin = 3; // reset is attached to digital pin 3\nint spectrumValue=0; // to hold a2d values\n\nconst int numReadings2 = 20;\nconst int numReadings3 = 20;\n\nint index2 = 0;                  // the index of the current reading\nint index3 = 0;\nfloat total2 = 0.0;                  // the running total\nfloat total3 = 0.0;\nint average = 0;                // the average\nfloat cycles2 = 0.0;\nfloat cycles3 = 0.0;\nfloat cyclesS2 = 0.0;\nfloat cyclesS3 = 0.0;\nfloat cyclesV = 0.0;\nfloat cyclesV3 = 0.0;\nfloat sumBPM2 = 0.0;\nfloat sumBPM3 = 0.0;\n//  int BperM2 = 0;\nfloat spectrumSum2 = 0.0;\nfloat spectrumSum3 = 0.0;\nfloat average2 = 0.0; //using\nfloat average3 = 0.0;\n//  int average4 = 0;\n//  int average5 = 0;  // not using these, I was thinking it would be cool to work in the different frequencies later.\n//  int average6 = 0;\n//  int average7 = 0;\nint holdValue1 = 0;   // this is for later if I want to do more with different frequencies.    frequency of 63 Hz\nint holdValue2 = 0;   // I'm currently using this one which is assigned the frequency of 160 Hz\nint holdValue3 = 0;   // not using frequency of 400 Hz\nint holdValue4 = 0;   // not using this one either    frequency of 1 KHz\nint holdValue5 = 0;   // or this    frequency of 2.5 KHz\nint holdValue6 = 0;   //or this    frequency of 6.25 KHz\nint holdValue7 = 0;   // or this frequency of 16 Hz (seems more like 15 KHz)\nint averageBPM1 = 0;  // not using\nfloat aveBPM2 = 0.0;  // I'm using this\nfloat aveBPM3 = 0.0;  // and this\nint holdBeat2 = 0;\nint holdBeat3 = 0;\nint beat2 = 0;\nint beat3 = 0;\nunsigned long lastCycle ;\nint saveBPM2 = 0;\nunsigned long saveAverage = 0;\nint b = 0;\nunsigned long timesAround = 0;\nfloat aveBPerM = 0.0;\nfloat aveBPerM3 = 0.0;\nint Fruit = 0;\nint Fruit3 = 0;\nfloat variance = 0;\nfloat variance3 = 0;\nfloat C = 0;\nfloat C3 = 0;\nfloat averageVariance = 0;\nfloat averageVariance3 = 0;\nfloat greatestAverageVariance;\nfloat greatestVariance = 0;\nunsigned long varianceSum = 0;\nunsigned long varianceSum3 = 0;\nint holdBPMDelay = 0;\ndouble BPMDelay = 0.0;\nint trigger2 = 0;\nint trigger3 = 0;\nint greg2 = 0;\nint greg3 = 0;\nint byteSending = 1;\n//int toTransfer = 32767;\nint Shift = BPMDelay;\nint mask = 0xFF;\nunsigned char toSend = 0;\nint twenty2 = 0;\nint twenty3 = 0;\n\nvoid setup()\n{\nSerial.begin(9600);\nWire.begin(0x52); // Start I2C Bus as slave\nWire.onRequest(requestInt);\nSerial.print(\"new Session\");\npinMode(analogPin, INPUT);\npinMode(strobePin, OUTPUT);\npinMode(resetPin, OUTPUT);\nanalogReference(DEFAULT);\ndigitalWrite(resetPin, LOW);\ndigitalWrite(strobePin, HIGH);\ndelayMicroseconds(36);\ntimer.setInterval(1250, BPM2);\ntimer.setInterval(1250, BPM3);\ntimer.setInterval(1250, ResetSec);\n//  timer.setInterval(11250, AverageBPM2);\n//  timer.setInterval(11250, AverageBPM3);\n//  timer.setInterval(11250, ResetTenSec);\ntimer.setInterval(11250, Fruity); // ran out of ideas for names so I chose Fruity all it does is switch a boolean statement\n{\n}\n{\n}\n}\n\ndreggory", null, "#2\nJan 14, 2013, 05:55 am\nCode: [Select]\nvoid requestInt()\n{\ndelay(8);\nif (byteSending == 1) //send packet 1\n{\ntoSend = Shift & mask;\nShift = Shift >> 8;\nWire.send(toSend);\nbyteSending = 2;\n}\nelse if (byteSending == 2) //send packet 2\n{\ntoSend = Shift & mask;\nShift = Shift >> 8;\nWire.send(toSend);\nbyteSending = 3;\n}\nelse if (byteSending == 3) //send packet 3\n{\ntoSend = Shift & mask;\nShift = Shift >> 8;\nWire.send(toSend);\nbyteSending = 4;\n}\nelse if (byteSending == 4) //send packet 4\n{\ntoSend = Shift & mask;\nShift = Shift >> 8;\nWire.send(toSend);\nbyteSending = 1;\n//initialization for next turn\nShift = BPMDelay;\ntoSend = 0;\n}\n\n}\n\nint Strobe()\n{\ndigitalWrite(resetPin, HIGH);\ndelay(33);\ndigitalWrite(resetPin, LOW);\n\nfor(int b=0; b < 7; b++)\n{\ndigitalWrite(strobePin, HIGH);\ndelayMicroseconds(18);\nif (b==0)\n{\ndigitalWrite(strobePin, LOW);\ndelayMicroseconds(36);\n//                    Serial.print(\" s1=\");\n//                    Serial.print(spectrumValue);\n//                    Spacing();\nholdValue1 = spectrumValue;\n\n}\nif (b==1)\n{\ndigitalWrite(strobePin, LOW);\ndelayMicroseconds(36);\n//                    Serial.print(\" s2=\");\n//                    Serial.print(spectrumValue);\n//                    Spacing();\nholdValue2 = spectrumValue;\n\n}\nif (b==2)\n{\ndigitalWrite(strobePin, LOW);\ndelayMicroseconds(36);\n//                    Serial.print(\" s3=\");\n//                    Serial.print(spectrumValue);\n//                    Spacing();\nholdValue3 = spectrumValue;\n\n}\nif (b==3)\n{\ndigitalWrite(strobePin, LOW);\ndelayMicroseconds(36);\n//                    Serial.print(\" s4=\");\n//                    Serial.print(spectrumValue);\n//                    Spacing();\nholdValue4 = spectrumValue;\n\n}\nif (b==4)\n{\ndigitalWrite(strobePin, LOW);\ndelayMicroseconds(36);\n//                    Serial.print(\" s5=\");\n//                    Serial.print(spectrumValue);\n//                    Spacing();\nholdValue5 = spectrumValue;\n\n}\nif (b==5)\n{\ndigitalWrite(strobePin, LOW);\ndelayMicroseconds(36);\n//                    Serial.print(\" s6=\");\n//                    Serial.print(spectrumValue);\n//                    Spacing();\nholdValue6 = spectrumValue;\n\n}\nif (b==6)\n{\ndigitalWrite(strobePin, LOW);\ndelayMicroseconds(36);\n//                    Serial.print(\" s7=\");\n//                    Serial.print(spectrumValue);\n//                    Spacing();\nholdValue7 = spectrumValue;\n\n}\n\n}\n//              Serial.println();\n\n}\n\nint Spacing()\n{\nif (spectrumValue < 10)\n{\nSerial.print(\"    \");\n}\nelse if (spectrumValue < 100 )\n{\nSerial.print(\"   \");\n}\nelse\n{\nSerial.print(\"  \");\n}\n}\n\nint AverageSpectrum2And3()\n\ncyclesS2++;\ncyclesS3++;\n\nspectrumSum2 += holdValue2;\nspectrumSum3 += holdValue3;\naverage2 = ((spectrumSum2)/(cyclesS2));\naverage3 = ((spectrumSum3)/(cyclesS3));\n//  Serial.print(\" spectrum ave1=\");\n//  Serial.print(average2);\n//  Serial.print(\" add it up=\");\n//  Serial.print(cycles1);\n//  Serial.print(\" Spectrum Sum =\");\n//  Serial.print(spectrumSum1);\n//  if (cycles1 - lastCycle > 4)    //count to four make a new line\n// {\n//\n//   Serial.println();                //  make sure to un-comment\n//   lastCycle = cycles1;\n//\n// }\n}\n\nvoid DynamicC()\n{\ncyclesV++;\nif (holdValue2 > average2)\n{\nvariance = ((holdValue2)-(average2));\n}\nif (average2 > holdValue2)\n{\nvariance = ((average2)-(holdValue2));\n}\n//  Serial.print(\" variance =\");\n//  Serial.print(variance);\n//  if (variance > greatestVariance)\n//       {\n//         greatestVariance = variance;\n//         Serial.print(\" GreatestVariance=\");\n//         Serial.print(greatestVariance);\n//       }\n\nvarianceSum += variance;\n//  Serial.print(\" varianceSum =\");\n//  Serial.print(varianceSum);\n//  Serial.println();\naverageVariance = ((varianceSum)/(cyclesV));\n//  Serial.print(\" aveVariance =\");\n//  Serial.print(averageVariance);\n//  Serial.println();\n//  if (averageVariance > greatestAverageVariance)\n//  {\n//    greatestAverageVariance = averageVariance;\n//    Serial.print(\" greatest Average Variance =\");\n//    Serial.print(greatestAverageVariance);\n//    Serial.println();\n//  }\n\nC = (((0.00125)*(averageVariance))+(1));\n//  Serial.print(\" C =\");\n//  Serial.print(C);\n//  Serial.println();\n}\n\nvoid DynamicC3()\n{\ncyclesV3++;\nif (holdValue3 > average3)\n{\nvariance3 = ((holdValue3)-(average3));\n}\nif (average3 > holdValue3)\n{\nvariance3 = ((average3)-(holdValue3));\n}\n//  Serial.print(\" variance =\");\n//  Serial.print(variance);\n//  if (variance > greatestVariance)\n//       {\n//         greatestVariance = variance;\n//         Serial.print(\" GreatestVariance=\");\n//         Serial.print(greatestVariance);\n//       }\n\nvarianceSum3 += variance3;\n//  Serial.print(\" varianceSum =\");\n//  Serial.print(varianceSum);\n//  Serial.println();\naverageVariance3 = ((varianceSum3)/(cyclesV3));\n//  Serial.print(\" aveVariance =\");\n//  Serial.print(averageVariance);\n//  Serial.println();\n//  if (averageVariance > greatestAverageVariance)\n//  {\n//    greatestAverageVariance = averageVariance;\n//    Serial.print(\" greatest Average Variance =\");\n//    Serial.print(greatestAverageVariance);\n//    Serial.println();\n//  }\n\nC3 = (((0.00125)*(averageVariance3))+(1));\n//  Serial.print(\" C =\");\n//  Serial.print(C);\n//  Serial.println();\n}\n\ndreggory", null, "#3\nJan 14, 2013, 05:56 am\nCode: [Select]\nvoid Beat2()\n{\nif(trigger2 == 0)\n{\nif((holdValue2) > ((average2)*(C)))\n{\n\nbeat2++;\n\ntrigger2 = 1;\n}\n}\nif((holdValue2) < ((average2)*(C)))\n{\ntrigger2 = 0;\n}\n\n}\n\nvoid Beat3()\n{\nif(trigger3 == 0)\n{\nif((holdValue3) > ((average3)*(C3)))\n{\n\nbeat3++;\n\ntrigger3 = 1;\n}\n}\nif((holdValue3) < ((average3)*(C3)))\n{\ntrigger3 = 0;\n}\n\n}\n\nvoid BPM2()\n{\n//  // subtract the last reading:\n//  total2 -= readings2[index2];\n// read the beats\n// add the reading to the total:\nfor (int i2 = 0; i2 < 19; ++i2)\n{\n\n}\nSerial.print(\" TOTAL 2=\");\nSerial.print(total2);\nSerial.println();\n// advance to the next position in the array:\nindex2++;\n\n// if at the end of the array...\nif (index2 >= numReadings2)\n\n// ...wrap around to the beginning:\nindex2 = 0;\n}\n\n//  holdBeat2 = beat2;\n\nSerial.print(\" **** BPS 2 ****=\");\nSerial.print(beat2);\nSerial.println();\n\ngreg2 = 2;\n\n}\n\nvoid BPM3()\n{\n\n//  // subtract the last reading:\n//  total3 -= readings3[index3];\n\n// read the beats\n// add the reading to the total:\nfor (int i3 = 0; i3 < 19; ++i3)\n{\n//      if(i3 == 0)\n//      {\n//        Serial.print(\"yaaaaaaaaaaaaahhhoooooooooo\");\n//      }\n}\nSerial.print(\" TOTAL 3=\");\nSerial.print(total3);\nSerial.println();\n\n// advance to the next position in the array:\nindex3++;\n\n// if at the end of the array...\nif (index3 >= numReadings3)\n\n// ...wrap around to the beginning:\nindex3 = 0;\n}\n\n//  holdBeat3 = beat3;\n\nSerial.print(\" **** BPS 3 ****=\");\nSerial.print(beat3);\nSerial.println();\n\ngreg3 = 2;\n\n}\n\nvoid AverageBPM2()\n{\nif(cycles2 < 20)\n{\ncycles2++;\n\n}\nSerial.print(\"cycles2=\");\nSerial.print(cycles2);\nSerial.println();\n\naveBPM2 = ((total2)/(cycles2));\naveBPerM = (aveBPM2 * 48);\ndelay(100);\nif (aveBPerM < 30)\n{\naveBPerM = 30;\n}\nif (aveBPerM > 500)\n{\naveBPerM = 500;\n}\n\ntotal2 = 0;\n\ngreg2 = 0;\n\n}\n\ndreggory", null, "#4\nJan 14, 2013, 05:57 am\nCode: [Select]\nvoid AverageBPM3()\n{\n\nif(cycles3 < 20)\n{\ncycles3++;\n}\nSerial.print(\"cycles3=\");\nSerial.print(cycles3);\nSerial.println();\naveBPM3 = ((total3)/(cycles3));\naveBPerM3 = (aveBPM3 * 48);\ndelay(100);\nif (aveBPerM3 < 30)\n{\naveBPerM3 = 30;\n}\nif (aveBPerM3 > 500)\n{\naveBPerM3 = 500;\n}\n\ntotal3 = 0;\n\ngreg3 = 0;\n\n}\n\nvoid ResetSec()  //ran every second\n{\n// delay(10);\nlastCycle = 0;\n//  holdBeat2 = beat2;\n//  holdBeat3 = beat3;\nbeat2 = 0;\nbeat3 = 0;\nvariance = 0;\nvariance3 = 0;\nvarianceSum = 0;\nvarianceSum3 = 0;\naverageVariance = 0;\naverageVariance3 = 0;\ncyclesV = 0;\ncyclesV3 = 0;\n\n}\n\nvoid ResetTenSec2()  //was used to run every ten seconds\n{\ndelay(15);\nspectrumSum2 = 0;\n\ncycles2 = 0;\n\nsumBPM2 = 0;\n\ncyclesS2 = 0;\n\n{\n}\n\n}\n\nvoid ResetTenSec3()  //was used to run every ten seconds\n{\ndelay(15);\n\nspectrumSum3 = 0;\n\ncycles3 = 0;\n\nsumBPM3 = 0;\ncyclesS3 = 0;\n\n{\n}\n\n}\n\nvoid Fruity()\n{\ndelay(20);\nFruit3 = 0;\n}\n\nvoid FruityOs()\n{\nFruit = 0;\n}\n\nvoid software_Reset() // Restarts program from beginning but does not reset the peripherals and registers\n{\nasm volatile (\"  jmp 0\");\n}\n\nvoid loop()\n{\n\nStrobe();\nif (holdValue2 > 130)\n{\nFruit = 0;\ndelay(10);\ndo\n{\nStrobe();\nAverageSpectrum2And3();\nDynamicC();\nBeat2();\ntimer.run();\nif (greg2 == 2)\n{\nAverageBPM2();\nif (aveBPerM > 35)\n{\nBPMDelay = (((60)/(aveBPerM))*(1000));\ndelay(50);\n//             Serial.print(\"^^^^BPMDelay in 2^^^^=\");\n//             Serial.print(BPMDelay);\n//             Serial.print(\"****BPMDelay is set in 2****\");\n}\n\nSerial.print(\"      ******* Average BPM 2 *******= \");\nSerial.print(aveBPerM);\nSerial.print(\"      BPMDelay =\");\nSerial.print(BPMDelay);\nSerial.println();\n}\n\nif ((holdValue2< 100) || (aveBPerM == 30))\n{\nResetTenSec2();\nSerial.print(\"holdValue2 =\");\nSerial.print(holdValue2);\nFruit = 2;\nFruit3 = 2;\n}\n\n}while(Fruit == 0);\n}\n\nif ((holdValue3 > 130)&&(Fruit3 == 2))         //detect third spectrum if fail to run spectrum 2\n{\n\ndelay(10);\ndo\n{\nStrobe();\nAverageSpectrum2And3();\nDynamicC3();\nBeat3();\ntimer.run();\nif(greg3 == 2)\n{\nAverageBPM3();\nif (aveBPerM3 > 35)\n{\nBPMDelay = (((60)/(aveBPerM3))*(1000));\ndelay(50);\nSerial.print(\"BPMDelay is set in 3\");\n}\n\nSerial.print(\"      ******* Average BPM 3 *******= \");\nSerial.print(aveBPerM3);\nSerial.print(\"      BPMDelay =\");\nSerial.print(BPMDelay);\nSerial.println();\n\n}\n\nif ((holdValue3< 100) || (aveBPerM3 == 30))\n{\nResetTenSec3();\n\nFruit3 = 0;\n}\n\n}while(Fruit3 == 2);\n}\n\nelse\n{\nSerial.print(\"Waiting...\");\nSerial.println();\nBPMDelay = 115;\n}\n\n//\n//if(Fruit == 2)\n//{\n//  if (aveBPerM < 35)\n//    {\n//      Fruit = 0;\n//    }\n//  delay(100);\n//  ResetSec();\n////  ResetTenSec();\n//}\n\n}\n\ndreggory", null, "#5\nJan 14, 2013, 05:58 am\nI wish I could have posted it all in one but I had to cut it up, it's a large sketch.\n\nReferences:\nhttp://www-scf.usc.edu/~ise575/a/projects/mooser/BeatDetector.htm\n\nhmm you know... as I was checking the bookmarks I had used for reference, I noticed this ^ was the only link that wasn't broken. expired web pages sorry.\n\nmarco_c", null, "#6\nJan 14, 2013, 08:14 am\nYou can attach files to a posting by clicking on the Additional option... At the bottom of the reply/post screen.\nArduino Libraries https://github.com/MajicDesigns?tab=Repositories\nParola for Arduino https://github.com/MajicDesigns/Parola\nArduino++ blog https://arduinoplusplus.wordpress.com\n\ndreggory", null, "#7\nJan 14, 2013, 09:03 am\nAwesome, thanks!\n\nok here is my sketch in one file\n\nMagician", null, "#8\nJan 15, 2013, 01:28 am\nError 503 Service Unavailable\n\nService Unavailable\nGuru Meditation:\n\nXID: 1087605664\n\nVarnish cache server\n\nWhat does it mean???\n\ndreggory", null, "#9\nJan 15, 2013, 08:08 pmLast Edit: Jul 20, 2014, 03:26 am by dreggory Reason: 1\nI don't know. Seems it didn't upload the sketch right. I'll upload it again.\n\nEDIT: the link started working and now both links work", null, "dreggory", null, "#10\nJan 27, 2013, 06:15 pm\n\nHere's a video of the whole project put together\n\nzmanpac", null, "#11\nJul 12, 2014, 05:13 pm\nI am new to the Arduino Blogs. I converted from Picaxe to Ardunio since they are much more versatile.  I was hoping you have the setup for the RGB posted some where I can look at how you setup the RGB display. You had mentioned it is a Dual Ardunio configuration on generates the BPM and the other displays the graphics. Was wondering what you use the output from ardunio 1 and inputs to the arduino 2 and the pinouts or if you have the code and schematic?\n\nVery Much appreciate and I think you have done an awesome job.\n\ndreggory", null, "#12\nJul 13, 2014, 12:12 amLast Edit: Jul 13, 2014, 12:14 am by dreggory Reason: 1\nI will share the whole project with you, no problem.\njust one thing... the animation program I used, originated from someone Else's code and I can't remember who I got it from, I adapted it heavily for my use but I still feel it would be rude/plagiarism to not cite who wrote it originally. so instead of posting it for all to see, I will share it over PM. my schematics are terrible, I made them in fritzing and only worked on cleaning up the PCB version(not the schematic.)\nlater I wanted to use 256 LEDs so I added 4 more shift registers(so the video that is posted is only 64 of the 256 LEDs)\nit could be sorted through and cleaned up, but I am too busy to get involved in this.\nI can send you the files and all the info you will need; but I can't help much more than that, sorry.\n\nzmanpac", null, "#13\nJul 19, 2014, 09:03 pm\nThat is fine I was hoping you had something where if I came up with some or a version different I be glad to share it you and all in here. I appreciate anything you can provide even if the schematics are done with pencil and xerox and scanned in is fine so long as they are legible I can create the schematic using DesignSpark tool for schematics.\n\nBy the way I make PCB using my CNC machine. If a PCB needs to be proto let me hace the schematic I can cut it on the CNC machine.\n\ndreggory", null, "#14\nJul 20, 2014, 03:10 am\nI'll just post it for all, if the person who wrote the code recognizes it I will just add the recognition.\nattached are the files needed\n\nalso the animations are written for common cathode LEDs, so if your matrix is common anode, you will have to invert the binary numbers (example: 01000010 will be rewritten as 10111101, and change the methods that look like this:\nCode: [Select]\n\nvoid allRed() //\n// turns on all RED\n{\ndigitalWrite(latchpin, LOW);\nshiftOut(datapin, clockpin, MSBFIRST, 255);\nshiftOut(datapin, clockpin, MSBFIRST, 0);\nshiftOut(datapin, clockpin, MSBFIRST, 0);\nshiftOut(datapin, clockpin, MSBFIRST, 255);// cathodes\ndigitalWrite(latchpin, HIGH);\n}\n\nneed to be changed to this:\n\nvoid allRed() //\n// turns on all RED\n{\ndigitalWrite(latchpin, LOW);\nshiftOut(datapin, clockpin, MSBFIRST, 0);\nshiftOut(datapin, clockpin, MSBFIRST, 255);\nshiftOut(datapin, clockpin, MSBFIRST, 255);\nshiftOut(datapin, clockpin, MSBFIRST, 0);// cathodes\n\ndigitalWrite(latchpin, HIGH);\n}\n\nBut you can ignore the code changes if you are using a common cathode matrix.\n\nIMPORTANT!\nin order for the whole thing to work you need to use the arduino IDE version 22 or it will not compile. (I didn't feel like porting it over to the new IDE)\n\nand lastly the breadboard/PCB project was made in fritzing which is a free PCB prototyping software.\n\nGo Up" ]
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https://curriculum.illustrativemathematics.org/HS/students/2/4/9/index.html
[ "# Lesson 9\n\nUsing Trigonometric Ratios to Find Angles\n\n• Let’s work backwards to find angles in right triangles.\n\n### 9.1: Once More with the Table\n\nA triangle with side lengths 3, 4, and 5 is a right triangle by the converse of the Pythagorean Theorem. What are the measures of the acute angles?\n\n### 9.2: From Ratios to Angles\n\nFind all missing side and angle measures.\n\nA good rule of thumb for a safe angle to use when leaning a ladder is the angle formed by your body when you stand on the ground and hold your arms out parallel to the ground.\n\n1. What are the angles in the triangle formed by your body and the ladder?\n2. What are the angles in the triangle formed by the ladder, the ground, and the railing? Explain or show your reasoning.\n3. You have a 13 foot long ladder and need to climb to a 12 foot tall roof.\n1. If you put the top of the ladder at the top of the wall, what angle is formed between the ladder and the ground?\n2. Is it possible to adjust the ladder to a safe angle? If so, give someone instructions to do so. If not, explain why not.\n\nPeople have various proportions to their body. Suppose that someone’s height to arm ratio is $$5:1$$.\n\n1. What are the angles in the triangle formed by their body and the ladder?\n2. How far off is this from the $$4:1$$ safe angle?\n3. What could this person do to make the ladder closer to the safe ladder angle?\n\n### Summary\n\nUsing trigonometric ratios and a calculator, the missing sides and angles of right triangles can be found.\n\nUsing the right triangle table we can estimate angle measures as in previous lessons. However, with a calculator, we can find angles more precisely.\n\nThe side opposite angle $$A$$ is 3 units long, and the side adjacent to $$A$$ is 12 units long. So to find angle $$A$$, we write an equation using tangent: $$\\tan(\\alpha) = \\frac{3}{12}$$. To find the measure of angle $$A$$ we ask the calculator, “What angle has a tangent of $$\\frac{3}{12}$$?” To ask that, we use arctangent by writing $$\\arctan \\left(\\frac{3}{12} \\right)$$. If we know the cosine, we use arccosine to look up the angle, and if we know the sine, we use arcsine. So $$\\alpha=\\arctan \\left( \\frac{3}{12} \\right)$$, which means angle $$A$$ measures about 14 degrees.\n\nAngle $$B$$ can be calculated using another trigonometric equation or the Triangle Angle Sum Theorem. Let's use arctangent again. We know $$\\tan(\\theta)=\\frac{12}{3}$$, so $$\\theta=\\arctan \\left( \\frac{12}{3} \\right)$$, which is about 76 degrees. This matches the answer we get with the Triangle Angle Sum Theorem: $$180-90-14 = 76$$.\n\n### Glossary Entries\n\n• arccosine\n\nThe arccosine of a number between 0 and 1 is the acute angle whose cosine is that number.\n\n• arcsine\n\nThe arcsine of a number between 0 and 1 is the acute angle whose sine is that number.\n\n• arctangent\n\nThe arctangent of a positive number is the acute angle whose tangent is that number.\n\n• cosine\n\nThe cosine of an acute angle in a right triangle is the ratio (quotient) of the length of the adjacent leg to the length of the hypotenuse. In the diagram, $$\\cos(x)=\\frac{b}{c}$$.\n\nThe sine of an acute angle in a right triangle is the ratio (quotient) of the length of the opposite leg to the length of the hypotenuse. In the diagram, $$\\sin(x) = \\frac{a}{c}.$$\nThe tangent of an acute angle in a right triangle is the ratio (quotient) of the length of the opposite leg to the length of the adjacent leg. In the diagram, $$\\tan(x) = \\frac{a}{b}.$$" ]
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https://hackaday.io/project/24834-3rgb-image-lossless-compression-format/log/60183-recursive-definition-of-a-2d-signal
[ "# Recursive definition of a 2D signal\n\nA project log for 3RGB image lossless compression format\n\nTiled CODEC for still and moving pictures based on the 3R compaction algorithm\n\nEverybody today is familiar with the representation of a picture, as a two-dimensional array of discrete values. 3RGB achieves some compression by adding recursively higher orders of granularity, where groups of pixels are summarised by their range, as a pair of values that represent the minimal value of pixel(s) in the group, as well as the largest value (relative to the minimum value).\n\nThe whole picture, made of discrete and inherently bounded values (usually from 0 to 255), is condensed in several step down to a pair of {min, amplitude} values. This is a significant departure from the classic frequency analysis transforms (DCT, wavelets...) where the \"root value\" is usually a single DC component.\n\nCompression is possible because the amplitude is related to the minimum through the inherent boundaries of the input sample: minimum + amplitude max.bound\n\nIn other words, amplitude max.bound - minimum\n\nGiven these two informations about a block of samples, each sample can be represented with less bits, for example with a bitstream of truncated binary codes (also called phase-i or economy code) or with 3R (which has the added advantage of also compacting closely related value tighter).\n\nOnce the decoder receives the \"root\" pair {min,amp}, it can decode a first block that contain the array for the first/highest level of min values. This array can be encoded as a simple concatenation of phase-in codes or with 3R for better compaction. At the same time, the first decode stage can generate the boundaries for the following array of amplitudes. Each individual min value will implicitly limit the range of the respective amplitude.\n\nAt the lowest level, each picture tile will have its minimum and amplitude values provided by the next higher level of arrays, and only an array of amplitudes is encoded. The only special case concerns the delta pictures: the amplitude of the tile is multiplied by 2 (unless it's 0, which means \"nothing to encode here, just fill the tile with a constant color\") and the LSB is a flag that selects between absolute or relative coding (the values are relative to the precedent value of the tile).\n\n## Discussions", null, "" ]
[ null, "https://analytics.supplyframe.com/trackingservlet/impression", null ]
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https://chemistry.stackexchange.com/questions/100162/are-there-any-reducing-agents-for-reduction-of-nitrobenzene-to-aniline-other-tha/100164
[ "# Are there any reducing agents for reduction of nitrobenzene to aniline other than Sn/HCl?\n\nI was taught that Sn/HCl should be used for reduction of nitrobenzene to aniline. Can I ask if there are any alternative reducing agents (for example, LiAlH4?) to do the same task?\n\nHydrogenation of aromatic nitro groups (over, for example, $$\\ce{Pd/C}$$) usually results in reduction to the corresponding aniline and is probably easier to carry out and work up as compared to $$\\ce{Sn/HCl}$$.\nWikipedia has a list of several conditions for reduction of $$\\ce{ArNO2}$$ to $$\\ce{ArNH2}$$. It also states that $$\\ce{LiAlH4}$$ reduces it to the azo compound $$\\ce{ArN=NAr}$$, which is corroborated by March's Advanced Organic Chemistry (7th ed.).\nIn any case, from a practical perspective, you would probably want to avoid using $$\\ce{LiAlH4}$$ if there's a simpler alternative. It's dangerous and the workup is not fun." ]
[ null ]
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https://gateoverflow.in/308548/ullman-toc-edition-3-exercise-3-2-question-4-page-no-108
[ "Convert the following regular expressions to NFA's with $\\in-$transactions.\n1. $01^{*}$\n2. $(0+1)01$\n3. $00(0+1)^{*}$" ]
[ null ]
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https://www.hindawi.com/journals/amse/2014/189423/
[ "Research Article | Open Access\n\nVolume 2014 |Article ID 189423 | https://doi.org/10.1155/2014/189423\n\nXiaohu Zeng, Shifeng Wen, Mingxi Li, Gongnan Xie, \"Estimating Young’s Modulus of Materials by a New Three-Point Bending Method\", Advances in Materials Science and Engineering, vol. 2014, Article ID 189423, 9 pages, 2014. https://doi.org/10.1155/2014/189423\n\n# Estimating Young’s Modulus of Materials by a New Three-Point Bending Method\n\nRevised06 May 2014\nAccepted10 Jun 2014\nPublished17 Jul 2014\n\n#### Abstract\n\nA new test method based on the three-point bending test is put forward to measure Young’s modulus of materials. The simplified mechanical model is established to make theoretical derivation. This method has not only the advantages of simple specimen preparation and convenient loading device, but also higher precision than the traditional three-point bending method. The method is adopted to obtain Young’s modulus of the aluminum alloy 2024. The feasibility of the method has been demonstrated by comparisons with the corresponding results obtained from the finite element method and experiment method. And the influence of contact friction on the test accuracy is analyzed.\n\n#### 1. Introduction\n\nUltrahigh temperature materials have been used in many fields to protect structural components from the intrinsic high temperatures, for example, on cooled gas turbines blades or vanes and the vehicle nose-tip ; therefore how to obtain material properties accurately under the ultrahigh temperature environment is the urgent problem to be solved in current research field .\n\nThere are some methods to obtain the mechanical properties of ultrahigh temperature materials especially the film/substrate system (the thermal barrier coating material) at present, such as tensile test, indentation test, and tympanic membrane test. However, all these methods are limited by the temperature and sample processing. The tensile test has a certain requirements of number and size on the sample, and the requirements cannot be met in some cases. First, it invariably runs into problems such as fixture oxidation, strength reduction, and processing difficulties due to high temperature in the fixture design process. Second, the thermal barrier coating (TBC) systems cannot be processed into standard samples. Third, the method is not suitable for ultrahigh temperature environment [7, 8]. Therefore, the indentation test has been applied to some extent. The indentation test has some advantages. First, it is local and nondestructive. Second, the test data can be measured directly in the stress concentration location. Zorzi and Perottoni studied Young’s modulus by indentation test in 2013 . However, this method still has some problems. For example, its precision in ultrahigh temperature environment is hard to control because of the influence of thermal drift produced in the ultrahigh temperature testing process . And the mechanical properties of materials obtained under high pressure stress do not have a good agreement with the results under tensile stresses. The bulge test is designed to determine the mechanical properties of thin film/substrate materials, but there is no report on application of the tympanic membrane test in ultrahigh temperature environment.\n\nIn the paper, a new test method is put forward to make up for lack of the above tests, namely, simple bending test method. The method has not only the advantages of simple specimen preparation and convenient loading device of the traditional three-point bending test, but also higher precision because the location of maximum stress and maximum deflection is not in the same place. For the traditional three-point and four-point bending test, the bending deflection and failure locations of the sample both occur in the middle position, as well as the contact position of indenter and measuring deflection data position, so the sample’s bending deformation will greatly influence the accuracy of deflection . The emphasis of the paper is to clarify the basic theory of adopting the new method to determine the Young’s modulus of single phase materials. The Young’s modulus is a basic mechanical parameter to characterize elastic deformation properties of solid materials and also an important basis for selecting mechanical components in engineering design. The test principle is described in the following part. And the feasibility of the method has been demonstrated by comparisons with the corresponding results obtained from the finite element method and experiment method. It can provide a theoretical foundation for applying the method to determine the elastic modulus of ultrahigh temperature materials.\n\n#### 2. Description of the Test Principle\n\nThe simplified geometry model of new test device is shown in Figure 1. Taking the point as the origin of coordinates, the coordinate system is established.\n\nThe support reactions are obtained by equilibrium conditions and ; that is, where and are the support reactions of and , respectively, is the applying load, is the distance between two supports, and is the distance between the loading point and the support . Then the shear forces and bending moment equations of the sample can be obtained, as follows: where and are the shear forces of and , respectively, and and are bending moments of and , respectively. From the theory of Mechanics of Materials , the equation of bending deformation is obtained: where is the curvature radius of deflection curve, is the Young’s modulus of the sample, and is the inertia moment of -axis. Because the cross section height is far smaller than the span, the influence on deformation of the shear stress is very small and can be ignored. Equation (3) can be rewritten as\n\nFrom the theory of Higher Mathematics, the following equation can be obtained:\n\nThe deflection curve is nearly flat in the elastic deformation stage, so the value of is very small. The following equation is obtained:\n\nThen the following approximate differential equations of the deflection curve can be obtained from (2) and (7): where and are the deflections of and , respectively and , , , and are unknown coefficients. Consider\n\nThe boundary conditions and continuity and smoothness conditions are described as follows: Substituting (15) into (9), (10), (12), and (13), the following equations are obtained:\n\nAs can be seen from the above equation, the bending deflection of is proportional to the applying load. From (19), the following equation is obtained:\n\nTherefore, as long as the ratio of applying load and bending deflection in the elastic deformation stage is obtained, Young’s modulus of the sample can be calculated.\n\n#### 3. The Test Device and Materials\n\nThe new test system of determining Young’s modulus proposed in the paper is shown in Figure 2. The test system is refitted by INSTRON8871 electrohydraulic servo fatigue testing machine. The system is equipped with a temperature control system and a cooling system, and it can create a high temperature environment above 800°C. The displacement meter comes with INSTRON8871 electrohydraulic servo fatigue testing machine. Its accuracy is 0.1% which can meet the requirements of the test. The fatigue loading system of ISTRON8871 fatigue machine is adopted to apply loads.\n\nThe test device adopted in the paper is shown in Figure 3. One end of the above transitional connection rod 7 is connected directly to the electrohydraulic servo fatigue testing machine. And the other end is connected with indenter 9 by bolt connection. The two nuts 8 and 15 are used to fasten indenter 9 and the lower transitional connection rod 16. The fixing device 12 connected to the lower transitional connection rod 16 by bolt connection is used to fix sample 10. And there are removable grooves 14 below the fixing device. It can change the contact position of indenter and sample to ensure the adequacy of test results.\n\nThe material of sample is aluminum alloy 2024. Its Young’s modulus () is known by consulting document. And its composition and content are shown in Table 1 . The sample is in a rectangle shape with length of 50 mm, width of 20 mm, and height of 2 mm. The material of indenter is cast nickel-base superalloy K403. The deformation of the indenter relative to the bending deflection of the sample is too small, so its deformation is not taken into consideration during the test. As the influence of contact friction between the indenter and sample on bending deflection is very small, it is also ignored.\n\n Si Fe Cu Mn Mg Cr Zn Ti Others Al Single Total 0.50 0.50 3.8~4.9 0.30~0.9 1.2~1.8 0.10 0.25 0.15 0.05 0.15 allowance\n\n#### 4. The Finite Element Model\n\nAccording to the characteristics of the new test device, the finite element model can be simplified as shown in Figure 4. The finite element software ABAQUS is adopted to carry on the numerical simulation in the paper. In order to accurately simulate the stress and deformation gradient in the contact region, the grids of contact region are refined. As the deformation of the indenter relative to the bending deflection of the sample can be ignored, the indenter and supports are simplified as a rigid body without deformation in the finite element analysis in order to simplify the calculation. And the influence of friction is not taken into consideration in the process of numerical simulation. The research object in this paper is Young’s modulus of materials, so the ambient temperature is not taken into consideration temporarily. The boundary condition applied to the model is that the two supports are fixed, the indenter is limited to move along the direction only, and the initial state of the sample is in horizontality. The load is applied on the indenter.\n\n#### 5. Results and Discussions\n\nThe load spectrum adopted in numerical simulation is a monotonically increasing spectrum, as shown in Figure 5. Figure 6 shows the relationship between the bending deflection and the applying load of the reference point of the sample. From the figure we can see that the slope of the load spectrum has no effect on the numerical simulation results; that is to say, the load’s increase speed with time has no effect on the bending deflection of the sample. And the deformation of the sample has two obvious stages. The first stage is stage AB: the stage where the bending deflection of the reference point monotonically increases with the applying load, that is, the elastic deformation stage. The second stage is stage BC: the stage where the bending deflection of the reference point rapidly increases with the applying load, that is, the plastic deformation stage. From the second section we know that as long as the ratio of applying load and bending deflection in the elastic deformation stage is obtained, Young’s modulus of the sample can be calculated. Therefore, the focus of this paper is the elastic deformation stage, that is, stage AB as shown in the figure. Young’s modulus of the sample in finite element simulation is calculated. Consider . This result has a good agreement with the Young’s modulus obtained by consulting document.\n\n#### 6. Experiment Results and Discussions\n\nThe load spectrum adopted in the test is a monotonically increasing spectrum in the same as the numerical simulation. It is applied to the sample by the loading system of ISTRON8871 fatigue machine. The test data are set to record every 0.02 s. All the tests are carried out at room temperature. The above load spectrum is applied on the sample to test three times without difference. The relationship between the bending deflection and the applying load of the sample is shown in Figure 7.\n\nFrom Figure 7, we can see that the deformation of the sample has also two obvious stages: stage AB, the stage where the bending deflection of the reference point monotonically increases with the applying load, that is, the elastic deformation stage; stage BC, the stage where the bending deflection of the reference point rapidly increases with the applying load, that is, the plastic deformation stage. The focus of this paper is stage AB. The test data processing is shown in Table 2. The method adopted in calculation process is linear regression method. The test results show that Young’s modulus of materials obtained by the new test method proposed in the paper corresponds to the Young’s modulus obtained by consulting document well. And the feasibility of the method has been demonstrated.\n\n Sample number Young’s modulus 1 GPa 2 GPa 3 GPa GPa\n\n#### 7. Error Analysis\n\nThe test error is calculated as . The factors that produce the error may be that the influences of deformation of the indenter and contact friction are not taken into consideration. And the influence of contact friction is studied as follows.\n\nFigure 8 shows the relationship between the bending deflection and the applying load of the reference point of the sample under different friction forces by finite element simulation. The Young’s modulus is calculated, as shown in Table 3. From the figure, we can see that the contact friction has influence on determination of Young’s modulus and the bigger contact friction coefficient is, the greater Young’s modulus will be. The bigger contact friction coefficient will cause bigger shear stress to prevent the sample continually deforming, so the bending deflection will be smaller and Young’s modulus will be greater.\n\n Contact friction coefficient Young’s modulus 0 70.24 GPa 0.25 75.1 GPa 0.5 79.3 GPa\n\n#### 8. Concluding Remarks\n\nIn the paper, a new test method and test device based on a new three-point bending test are put forward to measure Young’s modulus of materials, and the main results are summarized as follows.(1)The new method has advantages of simple principle, convenient operation, and wide field of application. The feasibility of the test method and test device has been demonstrated through finite element simulation and test on aluminum alloy 2024.(2)The contact friction has influence on determination of Young’s modulus of materials and the bigger contact friction coefficient is, the greater Young’s modulus will be.\n\n#### Conflict of Interests\n\nThe authors declare that there is no conflict of interests regarding the publication of this paper.\n\n#### Acknowledgments\n\nThis research is supported by Provincial Natural Science Foundation research project of Shaanxi (2014JQ1005), Aeronautical Science Foundation of China (2012ZD53053), Aerospace Technology Support Fund of China (2013-HT-XGD), and Specialized Research Fund for the Doctoral Program of Higher Education of China (20126102120034).\n\n1. B. P. Bewlay, M. R. Jackson, J. C. Zhao, P. R. Subramanian, M. G. Mendiratta, and J. J. Lewandowski, “Ultrahigh-temperature Nb-silicide-based composites,” MRS Bulletin, vol. 28, no. 9, pp. 646–653, 2003. View at: Publisher Site | Google Scholar\n2. T. Billot, P. Villechaise, M. Jouiad, and J. Mendez, “Creep-fatigue behavior at high temperature of a UDIMET 720 nickel-base superalloy,” International Journal of Fatigue, vol. 32, no. 5, pp. 824–829, 2010. View at: Publisher Site | Google Scholar\n3. J. Codrington, P. Nguyen, S. Y. Ho, and A. Kotousov, “Induction heating apparatus for high temperature testing of thermo-mechanical properties,” Applied Thermal Engineering, vol. 29, no. 14-15, pp. 2783–2789, 2009. View at: Publisher Site | Google Scholar\n4. P. F. Giroux, F. Dalle, M. Sauzay et al., “Influence of strain rate on P92 microstructural stability during fatigue tests at high temperature,” Procedia Engineering, vol. 2, no. 1, pp. 2141–2150, 2010. View at: Google Scholar\n5. C. Courcier, V. Maurel, L. Rémy, S. Quilici, I. Rouzou, and A. Phelippeau, “Interfacial damage based life model for EB-PVD thermal barrier coating,” Surface and Coatings Technology, vol. 205, no. 13-14, pp. 3763–3773, 2011. View at: Publisher Site | Google Scholar\n6. Y. Furuya, K. Kobayashi, M. Hayakawa, M. Sakamoto, Y. Koizumi, and H. Harada, “High-temperature ultrasonic fatigue testing of single-crystal superalloys,” Materials Letters, vol. 69, pp. 1–3, 2012. View at: Publisher Site | Google Scholar\n7. R. L. Edwards, G. Coles, and W. N. Sharpe Jr., “Comparison of tensile and bulge tests for thin-film silicon nitride,” Experimental Mechanics, vol. 44, no. 1, pp. 49–54, 2004. View at: Publisher Site | Google Scholar\n8. O. M. Abdelhadi, L. Ladani, and J. Razmi, “Fracture toughness of bonds using interfacial stresses in four-point bending test,” Mechanics of Materials, vol. 43, no. 12, pp. 885–900, 2011. View at: Publisher Site | Google Scholar\n9. J. E. Zorzi and C. A. Perottoni, “Estimating Young's modulus and Poisson's ratio by instrumented indentation test,” Materials Science and Engineering A, vol. 574, pp. 25–30, 2013. View at: Publisher Site | Google Scholar\n10. J. M. Kranenburg, C. A. Tweedie, K. J. van Vliet, and U. S. Schubert, “Challenges and progress in high-throughput screening of polymer mechanical properties by indentation,” Advanced Materials, vol. 21, no. 35, pp. 3551–3561, 2009. View at: Publisher Site | Google Scholar\n11. H. Ivankovic, E. Tkalcec, R. Rein, and H. Schmidt, “Microstructure and high temperature 4-point bending creep of sol-gel derived mullite ceramics,” Journal of the European Ceramic Society, vol. 26, no. 9, pp. 1637–1646, 2006. View at: Publisher Site | Google Scholar\n12. F. P. Beer and E. R. Johnston Jr., Mechanics of Materials, McGraw-Hill, New York, NY, USA, 1992.\n13. W. F. Brown Jr., H. Mindin, and C. Y. Ho, Aerospace Structural Metals Handbook, vol. 3, CINDAS/Purdue University, Weat Latayette, IN, USA, 1994.\n14. HKS, ABAQUS User’s Manual, version 6.2." ]
[ null ]
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https://www.zora.uzh.ch/id/eprint/129624/
[ "", null, "# Observation of the $B_s^0 \\rightarrow J/\\psi \\phi \\phi$ decay\n\nLHCb Collaboration; Bernet, R; Müller, K; Steinkampf, O; Straumann, U; Vollhardt, A; et al (2016). Observation of the $B_s^0 \\rightarrow J/\\psi \\phi \\phi$ decay. Journal of High Energy Physics, 2016:40.\n\n## Abstract\n\nThe $B_s^0 \\rightarrow J/\\psi \\phi \\phi$ decay is observed in $pp$ collision data corresponding to an integrated luminosity of 3 fb$^{-1}$ recorded by the LHCb detector at centre-of-mass energies of 7 TeV and 8 TeV. This is the first observation of this decay channel, with a statistical significance of 15 standard deviations. The mass of the $B_s^0$ meson is measured to be $5367.08\\,\\pm \\,0.38\\,\\pm\\, 0.15$ MeV/c$^2$. The branching fraction ratio$\\mathcal{B}(B_s^0 \\rightarrow J/\\psi \\phi \\phi)/\\mathcal{B}(B_s^0 \\rightarrow J/\\psi \\phi)$ is measured to be $0.0115\\,\\pm\\, 0.0012\\, ^{+0.0005}_{-0.0009}$. In both cases, the first uncertainty is statistical and the second is systematic. No evidence for non-resonant $B_s^0 \\rightarrow J/\\psi \\phi K^+K^-$ or $B_s^0 \\rightarrow J/\\psi K^+ K^- K^+ K^-$ decays is found.\n\n## Abstract\n\nThe $B_s^0 \\rightarrow J/\\psi \\phi \\phi$ decay is observed in $pp$ collision data corresponding to an integrated luminosity of 3 fb$^{-1}$ recorded by the LHCb detector at centre-of-mass energies of 7 TeV and 8 TeV. This is the first observation of this decay channel, with a statistical significance of 15 standard deviations. The mass of the $B_s^0$ meson is measured to be $5367.08\\,\\pm \\,0.38\\,\\pm\\, 0.15$ MeV/c$^2$. The branching fraction ratio$\\mathcal{B}(B_s^0 \\rightarrow J/\\psi \\phi \\phi)/\\mathcal{B}(B_s^0 \\rightarrow J/\\psi \\phi)$ is measured to be $0.0115\\,\\pm\\, 0.0012\\, ^{+0.0005}_{-0.0009}$. In both cases, the first uncertainty is statistical and the second is systematic. No evidence for non-resonant $B_s^0 \\rightarrow J/\\psi \\phi K^+K^-$ or $B_s^0 \\rightarrow J/\\psi K^+ K^- K^+ K^-$ decays is found.\n\n## Statistics\n\n### Citations\n\nDimensions.ai Metrics\n5 citations in Web of Science®\n3 citations in Scopus®\n\n### Altmetrics\n\nDetailed statistics", null, "", null, "Licence:", null, "" ]
[ null, "https://www.zora.uzh.ch/images/uzh_logo_en.jpg", null, "https://www.zora.uzh.ch/129624/8.hassmallThumbnailVersion/observation.pdf", null, "https://www.zora.uzh.ch/129624/8.haspreviewThumbnailVersion/observation.pdf", null, "https://www.zora.uzh.ch/license_images/by-4.0-88x31.png", null ]
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https://testbook.com/question-answer/the-block-diagram-of-a-closed-loop-control-system--5d00df5afdb8bb63e85a6ba7
[ "", null, "# The block diagram of a closed-loop control system is shown in the figure. The values of k and kp are such that the system has a damping ratio of 0.8 and an undamped natural frequency ωn of 4 rad/s respectively. The value of kp will be __________ .", null, "Free Practice With Testbook Mock Tests\n\nThis question was previously asked in\n\nGATE IN 2017 Official Paper\n\n## Solution:\n\nWe know that characteristic equation.\n\n1 + G(s) H(s) = 0\n\n$$1 + \\frac{{k\\left( {1 + {k_P}s} \\right)}}{{s\\left( {s + 1} \\right)}} = 0$$\n\ns2 + s (1 + kKp) + k = 0\n\ns2 + 2ζωns + ω2n = 0\n\n$${\\omega _n} = \\sqrt k = 4$$\n\nk = 16\n\n2ζωn = 1+ kKp\n\n1 + kKp = 2 × 0.8 × 4 = 6.4\n\nkKp = 5.4\n\nkp = 0.337" ]
[ null, "https://cdn.testbook.com/qb_resources/desktop_pyp.jpg", null, "https://storage.googleapis.com/tb-img/production/19/06/GATE%20IN_2017_Official_Sunny%20_Nita_Aman_8Q_images_Nita_Q4.PNG", null ]
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https://www.drillingformulas.com/kick-tolerance-factor-calculation/
[ "# Kick Tolerance Factor Calculation\n\nThe following formula is the kick tolerance factor.", null, "Where;\n\nCasing shoe TVD in ft\n\nWell depth TVD in ft\n\nMaximum Allowable MW in ppg\n\nCurrent MW in ppg\n\nFor more understanding, let’s try to calculate the kick tolerance factor (KTF) using the following information:\n\nCasing shoe TVD = 5000 ft\n\nWell depth TVD = 11,000 ft\n\nMaximum Allowable MW = 12.9 ppg\n\nCurrent MW =9.8 ppg\n\nAccording to the equation above, you will be able to calculate like this;", null, "Kick Tolerance Factor =1.41 ppg\n\nPlease find the calculation sheet >>", null, "Kick Tolerance Factor Calculation\n\nShare the joy", null, "1.", null, "gas certificate says:" ]
[ null, "https://www.drillingformulas.com/wp-content/uploads/2011/01/Kick-Tolerance-Factor.jpg", null, "https://www.drillingformulas.com/wp-content/uploads/2011/01/Kick-Tolerance-Factor-Sample.jpg", null, "https://www.drillingformulas.com/wp-content/uploads/2010/04/excel_icon.jpg", null, "https://secure.gravatar.com/avatar/564cdfcbd738c8ab1493fcca6fbfe825", null, "https://secure.gravatar.com/avatar/e9d3ae963ad31377fb50374a828dbb57", null ]
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https://developer.mescius.com/spreadnet/docs/online-asp/FarPoint.Web.Spread~FarPoint.Web.Spread.SheetView~MoveColumn(Int32,Int32,Int32).html
[ " MoveColumn(Int32,Int32,Int32) Method | Spread ASP.NET 16\nColumn index at which to start the move (first column)\nNumber of columns to move\nDestination before which to copy the range of columns\nExample\n\nIn This Topic\nMoveColumn(Int32,Int32,Int32) Method\nIn This Topic\nMoves the column or columns to a specified location.\nSyntax\n```'Declaration\n\nPublic Sub MoveColumn( _\nByVal column As Integer, _\nByVal count As Integer, _\nByVal toColumn As Integer _\n) ```\n```'Usage\n\nDim instance As SheetView\nDim column As Integer\nDim count As Integer\nDim toColumn As Integer\n\ninstance.MoveColumn(column, count, toColumn)```\n```public void MoveColumn(\nint column,\nint count,\nint toColumn\n)```\n\n#### Parameters\n\ncolumn\nColumn index at which to start the move (first column)\ncount\nNumber of columns to move\ntoColumn\nDestination before which to copy the range of columns\nExample\nThis example uses the MoveColumn method." ]
[ null ]
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http://conspiracyoflight.com/SagnacRel/SagnacandRel.html
[ "", null, "The search for  new physics.\n\nThe Sagnac Effect: Does it Contradict Relativity?\n\nBy Doug Marett 2012\n\nA number of authors have suggested that the Sagnac effect contradicts the original postulates of Special Relativity, since the postulate of the constancy of the speed of light is violated in rotating systems. [1,2,3] Sagnac himself designed the experiment of 1913 to prove the existence of the aether . Other authors have attempted to explain the effect within the theoretical framework of relativity, even going as far as calling the effect “relativistic”. However, we seek in this paper to show how the Sagnac effect contravenes in principle the concept of the relativity of time and motion.  To understand how this happens and its wider implications, it is necessary to look first at the concept of absolute vs. relative motion, what the Sagnac effect implies about these motions, and what follows logically about the broader notions of space and time.\n\nEarly Ideas Regarding Rotation and Absolute Motion:\n\nThe fundamental distinction between translational and rotational motion was first pointed out by Sir Isaac Newton, and emphasized later by the work of Ernest Mach and Heinrich Hertz.  One of the first experiments on the subject was Newton’s bucket experiment of 1689. It was an experiment to demonstrate that true rotational motion cannot be defined as relative rotation of a body with respect to surrounding bodies – that true motion and rest should be defined relative to absolute space instead.   In Newton’s experiment, a bucket is filled with water and hung by a rope. If the rope is twisted around and around until it is tight and then released, the bucket begins to spin rapidly, not only with respect to the observers watching it, but also with respect to the water in the bucket, which at first doesn’t move and remains flat on its surface.", null, "Eventually, as the bucket continues to spin, the water starts to rotate as well, as can be seen from the concave shape of its surface. This concaving of the water shows that it is rotating, even though it is now at rest with respect to the co-rotating bucket. Newton pointed out it is not the relative motion of the bucket and water that causes the concavity of the water. The concaving of the water suggests that it is rotating with respect to something else, far more remote.  In Newton’s thinking, this showed rotation relative to say…absolute space. This is contrary to the idea that motions can only be relative.\n\nNow an observer who is sitting in the bucket and thereby rotating with it can see that the water is concave but may not be able to see or feel that he is rotating with respect to the surroundings, so the cause of the concavity may not be obvious, particularly if the rotation is slow and his view of the surroundings is blocked. On the other hand, an observer who is standing near the bucket (stationary) will see a concave shape of the water’s surface that is consistent with his observation of the bucket rotating in front of him. So although both observers see the effect of rotation (concave surface) they cannot necessarily agree on the cause.\n\nIt is sometimes argued that because a stationary frame can be identified (the observer near the bucket), there is no reason to ask “rotating with respect to what?” However, we wish to show that although multiple stationary frames can be identified depending on the circumstances, there is only one truly unique stationary frame, and this is the frame of absolute space, a.k.a “the fixed stars.”\n\nThe Sagnac Effect:\n\nWhile originally these ideas were applied to mechanical effects, with the Sagnac effect it becomes clear that the same principles apply to light, and the frame of reference in which light speed can be shown (or not shown) to be constant. Consider the Sagnac effect in the context of Newton’s bucket experiment. Instead of a rotating bucket, we have a rotating disk, and on the rim of the disk are two light paths originating from a source, one going clockwise and the other going counter-clockwise around the rim of the disk, and returning to a detector where their time of flight is measured. Since the transit time is usually very small, the differences in the return times are determined by a phase difference between the two returning beams. For the Sagnac effect, we call this Case 1. The speed of light (C) in the stationary frame is assumed to be C ~= 3 x 108 m/s.\n\nIf we now rotate the disk clockwise at velocity v, we find that the blue beam arrives back at the detector before the red beam – in fact, the difference in  the  velocity  of  light", null, "turns out to be 2*v, because the blue beam travels at C+v and the red beam at C-v.  Just like the concave water, all observers agree that the blue beam has arrived before the red beam, but again, like in the case of Newton’s bucket, the explanation of why depends on the perspective of the observer.  For the observer B it appears that the detector has approached the oncoming blue beam and has receded from the red beam, and it is the path length change with respect to the two beams that accounts for the apparent difference in the speed of light. On the other hand, the observer at A who is rotating with the disk has observed that the source and detector are stationary from his perspective, as are the paths, so he can only conclude (assuming no foreknowledge that he is rotating) that the light has actually travelled at two different velocities around the disk, being C+v and C-v.  Who is right? Well, the problem is this – from the point of view of special relativity, there are no preferred frames of reference, there are only relative frames of reference, so we can’t argue that frame B is better than frame A without defeating the underlying premise of the theory. In special relativity (SR), the position of the observer is supposed to be arbitrary. However, if we accept that perspective A is legitimate, we defeat the first postulate of SR, namely:\n\n“that light is always propagated in empty space with a definite velocity c which is independent of the state of motion of the emitting body”. – Einstein \n\nFurther, on the relativistic notion that ”it is impossible to detect motion by measuring differences in the speed of light”, as had seemingly been proven by the Michelson-Morley experiment, the Sagnac experiment shows that this can in fact be done. Also, observer A can determine that it is he who is rotating, and not his surroundings.\n\nNow to address the idea of stationary states we need to ask the following question: Is there a single unique observer (frame of reference) for which the speed of light is always constant in a Sagnac interferometer? In fact, the speed of light measured by observer B is not exactly constant either, a conundrum which we will explain below.", null, "In case 1 we recall that observer B in the lab was stationary and noted that the speed of light should be approximately C in both directions. Let’s expand our interferometer now to encompass the earth. We move our observer B to the North Pole, where he hovers, not rotating, and our observer A is placed in the lab on the earth’s surface, along with our source and detector. Our circular light path is now the equator of the earth. This configuration is in fact the experiment first proposed by Albert Michelson in 1904. Since the source, detector, and observer A are all on the surface of the earth, from their perspective they are stationary (excluding fore-knowledge). However, as before, A observes the speed of light reaching the detector after traveling around the globe to be C+v and C-v for the blue and red light paths, respectively, where v is the rotational velocity of the earth on its axis (465 m/s). Also, observer B at the North Pole can see the Earth rotating under him, so from his perspective the paths are again unequal and the speed of light is about constant from his perspective.\n\nSo now we can see that the lab frame is not a unique perspective – in case 1 the speed of light was approximately constant in the lab frame, in case 2 it is variable. However, even in case 1, the so called “stationary observer” in the lab will be able to detect the slow rotation of the much smaller lab-sized ring due to the earth spinning on its axis, of up to 15 degrees per hour, so the speed of light measured by this observer will not be constant either, but will differ by the much smaller amount C+/- (r * W) where r is the radius of the loop and W is the angular velocity of rotation at the observer’s latitude.\n\nHowever, the North Pole could be construed as an approximately stationary perspective for both case 1 and 2. So is the speed of light always constant with respect to the non-rotating Earth (the so called ECI frame used in GPS)? Again, we can show this not to be true.", null, "If we consider a third scenario (Case 3), where B is on the Sun and A is on the North Pole, and the light paths go around the orbit of the Earth, we find that now the North Pole perspective sees a variable speed for C (+/- the orbital velocity of the earth which is ~30 km/s) and only the observer on the sun sees an approximately constant C based on his observation that the two light beams have followed longer or shorter paths. But using the same reasoning as in case 2, the observer B on the sun will not measure exactly C either, but C+/- a much smaller value, again being C+/- (r * W) where r is the radius of the galaxy to our sun and W is the angular velocity of our sun around the galaxy.\n\nThis line of reasoning can continue to be expanded to larger and larger scales – for example, if A is standing on the Sun and B is at the center of the galaxy, then the center of the galaxy becomes the (approximate) preferred perspective. In fact, the only way that we can define an observer for which all circumstances the speed of light C is exactly constant in a circular path, is to place the observer B stationary with respect to the universe.", null, "Any light following a circular path with respect to B will always be observed to be exactly C, but only from this perspective. For all other velocity frames, the speed of light can be variable. In reaching this conclusion, we have established that the Sagnac interferometer implies that there must exist only one unique frame of reference for rotation, and this is the universe itself – absolute space, a.k.a. “the fixed stars”.  So to answer the question “with respect to what does the Sagnac interferometer rotate?” we must answer, “absolute space, the fixed stars.” But in doing so, we are compelled to reject the idea of the relativity of motion that is the cornerstone of Einstein’s theory.\n\nRelativists claim that the Sagnac effect can be accounted for by assuming a constant light speed on the rotating disk and by factoring in the time dilation of the rotating observer. This uses a backwards mathematical transformation from the stationary observer. This mathematical treatment, attributable to Langevin and repeated for example by E.J. Post, was decisively disproved by the Dufour and Prunier experiments. See Appendix A for the full discussion.\n\n______________________________________________________________________________________________\n\nGPS and the Sagnac Effect – GPS satellites must use clocks synchronized to the Earth Centered Inertial frame (ECI) – i.e. their satellite observer clock coincides with the non-rotating observer B on the North Pole of Case 2. If they use clocks synchronized to the observer A on the earth’s surface in the Earth Centered Earth Fixed frame (ECEF) of case 2 then the speed of light between the satellite and earth is no longer measured to be C, introducing a propagation range error and thereby an error in the receivers estimated position. This procedure has been used in practice now for decades.", null, "______________________________________________________________________________________________________________________\n\nThe Non-Relativity of Time\n\nThe Sagnac effect also does for time what it does for light. Many of us know Einstein’s famous twins paradox, where one twin stays on earth while the other travels in a rocket at near light speed, returning years later to meet his twin. When they meet again, the travelled twin is younger than the stay at home twin, because according to Einstein, the high velocity of the second twin meant that time progressed slower for him than the stay at home twin.\n\nThis thought experiment was criticized by Dingle, who aptly pointed out that if Einstein argues that a moving clock counts slower than a stationary one, but also presumes all motion to be relative, then we should be able to argue with equal justification that it is the earthbound twin who has moved away at velocity from the rocket bound twin, and that the former’s clock has counted slower instead.\n\nAn attempt to resolve this issue was the Hafele and Keating Experiment. Hafele and Keating [9,10] performed an experiment in 1971 to test Einstein's predictions regarding the dilation of time in clocks moved at some velocity with respect to a stationary clock. In the opening statement of the first of two papers on the subject, the authors refer to the debate surrounding the \"twins paradox\" and how an experiment with macroscopic clocks might provide an empirical resolution. In Hafele and Keating's experiment, they flew cesium clocks around the world in opposite directions near the equator, and then measured how much time they have gained or lost when they return to the start point as compared to a stationary cesium clock to which they were originally synchronized. The experiment of Hafele and Keating (HK) actually differs from the twins paradox in some important ways - firstly, in the paradox, the traveling twin moves in a straight line to a distant point then turns around and comes back, with a series of accelerations and decelerations. In the HK experiment, the journey is circular and the clocks never leave the earth’s influence. In following a circular path, the HK experiment becomes a Sagnac effect experiment where it is clocks that are moving rather than beams of light.", null, "From the perspective of the observer on the ground, the moving clocks of Hafele and Keating have experienced similar motions with respect to the stationary earth bound clock. If their motions were truly \"relative\", then by the rules of SR both should arrive back at the starting point having experienced the same amount of time dilation in the same direction.    The actual result was that one travelled clock increased in its time rate and the other decreased. This is the identical result to what would be expected for light in the Sagnac experiment Case 2, where one beam travels faster than expected, and the other slower than expected. And as before, it implies that the only observer who sees all clocks count at a rate proportional to their velocity is the observer at the North Pole, who is in the non-rotating, geocentric (ECI) frame of the earth.\n\nAs Murray pointed out, Hafele and Keating interpreted their result by following a revised form of relativity attributable to Builder, who they reference in their first paper in 1971. In Builder's \"Ether and Relativity\" he summarizes his view as follows :\n\n\"The relative retardation of clocks, predicted by the restricted theory of relativity, demands our recognition of the causal significance of absolute velocities. This demand is also implied by the relativistic equations of electrodynamics and even by the formulation of the restricted theory itself. The observable effects of absolute accelerations and of absolute velocities must be ascribed to interaction of bodies and physical systems with some absolute inertial system. We have no alternative but to identify this absolute system with the universe. Thus in the context of physics, absolute motion must be understood to mean motion relative to the universe, and any wider or more abstract interpretation of the \"absolute\" must be denied. Interactions of bodies and physical systems with the universe cannot be described in terms of Mach's hypothesis, since this is untenable. There is therefore no alternative to the ether hypothesis. \"\n\nBuilder’s conclusion that there must be \"an absolute inertial system\" re-introduces the preferred frame of Lorentz back into relativity, the very thing that Einstein sought to eliminate.\n\nIf we are to accept Builder's hypothesis of a universal preferred frame for motion (ether), then we are compelled to accept the other critical tenet of Lorentz over Einstein, that it is the speed of light, and not time, which is the fundamental variable in our universe. In accepting a notion of absolute motion, or a preferred frame, we are as Dingle says, rejecting the notion of \"relative motion\". Despite Einstein borrowing heavily from Lorentz's equations and theory, and the fact that both theories predict similar results for the same experiments, the two theories are philosophically antithetical - and as such the conclusions about the nature of space and time result in entirely different models of the universe.\n\nLet’s look at the implications of the Hafele and Keating Experiment on the relativity of time. Similar to our original case 3, we expand the travelling clock experiment to the orbit of the earth. Rather than planes, we put cesium clocks on two rockets and have them take off from the surface of the earth at equal escape velocities (say 12 km/s), but set to travel in opposite directions around the sun before returning to the earth, as in the figure below. The rockets maintain the same velocity with respect to the earth’s surface throughout their voyage. Both rockets arrive back on the earth at around the same time, and again at the same velocity with respect to the earth, at 12 km/s.  As expected, the rocket on the left going counter-clockwise arrives with its clock rate slowed with respect to the earth clock, and the rocket on the right travelling clockwise has had its clock rate sped up with respect to the earth clock, even though their velocity with respect to the earth’s surface clock was the same. From the perspective of the earth (assuming no foreknowledge that they are  rotating),  the  path  lengths  for  each  rocket", null, "appear to be the same. However, an observer B on the sun would see that the rockets had to travel at different speeds with respect to the sun in order to be travelling at the same speed with respect to the orbiting earth. The paths that they travelled would also appear to be different lengths when viewed from the sun. Only the observer on the sun would agree that the SR postulate that a clock decreases it rate with velocity is found to be (approximately) true. The earth observer A sees this to be false – for him, the clocks have travelled identical paths at identical speeds, yet the clocks have gone out of synchronization.  This invalidates the notion that the ECI frame of the earth is a unique frame of reference for rotation with which the laws of SR hold true.\n\nAgain, we can extrapolate this thought experiment to the edge of the galaxy, and further to the entire universe.  For clocks rotating around an observer, the only unique perspective where ultimately any rotational motion will always lead to a slowing of the clock rate is the stationary frame of the universe. In any other velocity frame, it will always be possible to create orbiting clocks that appear to speed up in rate with respect to our observer. This is the same conclusion as earlier for the speed of light, except now applied to time.\n\nThese rotation experiments have implied that there is a an absolute space with which we can define a reference speed of light C, and a preferred frame of reference for a time clock for which all other moving clocks will count at a slower rate.\n\nSince Einstein invented the idea of a variable time in order to explain the constancy of the speed of light, should it not follow then that if we have shown that the speed of light is not always constant for moving observers, then perhaps instead it is “time” that is constant, and the velocity of light is variable depending on our motion with respect to absolute space? Since translational experiments always show a constant measured speed of light, unlike the Sagnac experiment, perhaps apparent time dilation is our only indirect evidence that the speed of light is changed in a hidden way during translational motion. In order to prove this we would need to design a clock that doesn’t depend on the speed of light in its counting principle, and see if this clock is unaffected by time dilation. Again, we will turn to the Sagnac effect and rotation to design such a clock. But first, let’s examine the basis for time dilation in Einstein’s theory.\n\nThe Light Clock Analogy\n\nWhen we see a supernova explode in the sky, we do not say “that event just happened”, since we know the light from that event has taken perhaps thousands of years to reach us from some distant point in space. We would say, rather, that the explosion happened thousands of years ago. This is because our judgement of the simultaneity of events is independent of the speed of light. Rather the speed of light might be used to insinuate when an event could have happened in the past, based on the propagation delay. So why then would we believe that the rate of time itself depends on the speed of light? Consider Einstein’s analogy, often reproduced, of a light beam bouncing between two mirrors. In his analogy, two horizontal mirrors face each other and one mirror is spaced above the other by a distance d. A light pulse bounces vertically between the two mirrors as shown on the left.  The time it takes for the pulse of light to do a round trip (from the top mirror to the bottom and back) is twice the distance d divided by the speed of light. This is his light clock.", null, "Suppose the \"light clock\" were traveling sideways at a very high (but constant) speed of 0.5 C. Now the pulse would follow the \"saw tooth\" path shown on the right side of the drawing. Since light speed is the same as before (remember, the speed of light is not changed by the speed of its source), it will take longer to make a round trip. So our \"light clock\" takes longer to count out its intervals. Another way of saying this is that the clock \"ticks\" more slowly.\n\nBut here lies the problem. This analogy works because the speed of light is constant as seen by observer A. From his vantage point, he can see both clocks working, and although the speed of light is travelling at the same rate in scenario 2, it has taken a longer path to hit the moving mirrors. In taking a longer path, it has taken more time to reach the end (based on A’s clock – 1.1546 x longer or t/0.866 = t’). But because we have already said that the rate of time is independent of the speed of light, we cannot assume that “real” time has passed slower for observer B than observer A, we can only conclude that B’s clock has counted mechanically slower because it has taken longer for the light to travel the distance between scenario 2’s mirrors. Now, if observer B measures the speed of light in his moving frame using “real” time, i.e. a clock not affected by the speed of light, he should conclude that the speed of light is slower than C, namely 0.866*C = C’. This is exactly like the situation of the Sagnac interferometer – observer B is unaware that he is moving, so he perceives the mirrors to be opposite each other and separated by distance d, so he would measure the speed of light to be less than C. He is unaware that from the perspective of observer A, the path d has lengthened. So why does he conclude that the speed of light is instead C? This is because he measures the time elapsed using a clock that is based on the length of time it takes for light to pass between his own moving mirrors. In effect, his moving clock is calibrated by the speed of light seen in his own frame (which is 0.866*C), so consequently, his clock counts slower by exactly this amount (t/0.866). This makes him falsely measure that the light beam to have arrived at velocity C rather than C * 0.866.\n\nIf the speed of light was actually C in his frame, i.e. C = C’, then his moving light clock would count faster (now t = t’) then although he would still measure the speed of light to be C with his clock, the time rate difference between A and B would disappear! Since the slowing of clocks (apparent time dilation) does occur in reality, this implies, just like in the Sagnac experiment, that the speed of light is only actually C in the stationary frame, and in the moving frame it is actually C *  0.866, but is erroneously measured as being C in his frame due to the slowing error introduced into his moving clock by the longer light path.\n\nWe can go one step further – the light clock analogy explains why every translational experiment to measure the speed of light shows that the speed of light is always C – because the clocks used to measure C are calibrated to the speed of light in their frame, this built-in correction factor always adjusts the measured speed back to C. However, in every rotational experiment, it is trivial to measure a speed of light that is not C. This is because, as we showed with Hafale and Keating, clocks moving in rotating paths around a disk or sphere will go out of synchronization with each other – using a common origin as a reference and a circular path to the detector confounds the self-correcting ability of a light clock, and the underlying difference in the speed of light with respect to the origin is revealed.\n\nThe Sagnac Clock:\n\nAs we said earlier, in order to prove that time dilation is an illusion,  we would need to design a clock that does not depend on the speed of light in its counting principle, and see if this clock is unaffected by time dilation. Again, we turn to essential feature of the Sagnac effect, rotation, to design such a clock. This concept was first proposed in our article “The Paradox of the Clocks in the Canaries” in 2010.\n\nThe simplest time dilation effect to test is gravitational time dilation, the speeding up of clocks with altitude. According to Einstein, the rate of time follows the gravitational potential, so a clock on the surface of the earth should count more slowly than a clock at some high altitude, such as at the top of a mountain. Time is alleged to be moving into the future faster at the top of a mountain.  So we need a clock on the Earth that is unaffected by altitude. Let's say we use the rotation of the earth as our clock instead of a conventional cesium clock. After all, our human concept of time, and the units that compose it (days, hours, minutes, seconds) are already derived from the rotation period of the earth, so it should be the ideal clock for comparison. For each observer, we could use as our reference the position of the sun in the sky, or the position of a Foucault pendulum, or most preferably the position of the fixed stars in sidereal time. If we measure the sidereal day (23.93447 hours) as the time for one earth rotation relative to the vernal equinox, using a sighting star as a reference, we can then place an observer A at the top of a mountain, and another observer B at sea level. Preferably the mountain observer is exactly vertical over the sea level observer, and this vertical line passes through the center of the earth, as shown in the picture below. We use as our example the mountain Pico del Teide on Tenerife, with an altitude of 3718 meters. We give each observer A and B a cesium clock, and a sidereal clock.  The sidereal clocks mark a tick when the reference star is directly overhead each day. The cesium clocks mark a tick when 23.93477 hours has passed in their local time. We then compare the four clocks to determine if they are counting at the same rate.\n\nWhat happens is that when the reference star passes over the red vertical line passing through observer A and B, both sidereal clocks have no choice but to trigger at the same moment, differing only in the light propagation time between A and B. Because this propagation time is a fixed amount, it does not affect the rate at which the two clocks count – therefore the count rate is exactly the same. However, since both cesium clocks experience gravitational time dilation, the cesium clock at observer A must count faster than the cesium clock at observer B. If we had pre-synchronized the two clocks at sea level, then the cesium clock at observer A will begin to go out of synchronization with", null, "the other three clocks, and the amount of time gained by cesium clock A will continue to increase as the days pass. Clearly three of the 4 clocks agree on the amount of elapsed time – only cesium clock A reads in error. The fact that we have created clocks immune to time dilation implies that real time has never changed at the top of the mountain – apparent time dilation is an artifact of the difference in the speed of light at two difference points in space, and the use of clocks calibrated to the speed of light – i.e. the light clocks of Einstein.\n\nThis conclusion is also consistent with our reason – given the millions of years that the mountain has been around, the summit of Pico del Teide should by now be several minutes in the future compared to the land at its sea shore base. And yet they are physically connected. To believe that time dilation is real would require us to assume that walking up the mountain is the same as walking into the future - the land would have to form a bridge between two points in Einstein’s dimension of time.  It is somewhat trivial to travel to the top of the mountain and view through a telescope a clock tower in Puerto de la Cruz far below to verify that the top of the mountain and the sea side resort share the same present.\n\nThe author at the summit of Pico del Teide in 2012, checking the time. Puerto de la Cruz is clearly visible 3718 meters below. Still in the present", null, "Conclusions:\n\nIt is often argued that the predictions of Special and General Relativity have been continuously verified and that therefore the theory is unquestionable. However, other theories, such as Lorentz ether theories modified to take into account gravitational effects, can also make similar claims. There are in fact multiple mathematical routes by which a correct prediction can be arrived at, but these theories may imply very different interpretations of what our physical reality is. And this is at the heart of what is wrong with the theory of relativity – it may make successful predictions based on math, but implies a nature of time and space which are not only inconsistent with logic and  reason, but are even contradictory. Our interpretation is that these problems arise because of the switch from the absolute time / variable speed of light of Lorentz (1904) to relative time / absolute speed of light of Einstein (1905) – two stances which are only slightly different mathematically. In Lorentz’s favour, the Sagnac effect demonstrates that depending on the placement of the observer, it is possible to see this variable speed of light and to confound apparent time dilation. Our analysis of Einstein’s light clock has shown that it will always count in error if the speed of light is not C in the moving frame. And further, our thought experiment with sidereal clocks has shown that the entire premise that time dilation corresponds to a change in “real” time is highly questionable. So when countered with the argument that General Relativity can explain the Sagnac effect, we might ask, why bother? If time dilation is an illusion, then the entire 4D time-space continuum of Einstein should be considered, to use his own word for the aether, “superfluous.”\n\nReferences:\n\n Moon, Parry, Spencer, Domina Eberle, and Uma, Shama, (1991) The Sagnac Effect and the Postulates on the Velocity of Light.\n\n Ruyong Wang  (2003) Modifed Sagnac experiment for measuring travel-time difference between counter-propagating light beams in a uniformly moving fiber.\n\n Alexandre Dufour, Fernand Prunier  (1942) On the Fringe Movement Registered on a Platform in Uniform Motion\n\n Georges Sagnac (1913) Regarding the Proof for the Existence of a Luminiferous Ether using a Rotating Interferometer Experiment.\n\n Rizzi, Guido, (2003) The Relativistic Sagnac Effect: two derivations.\n\n A. Michelson, (1904) \"Relative Motion of Earth and Aether.\" Philosophical Magazine (6), 8, 716-719.\n\n Albert Einstein (1905) On the Electrodynamics of Moving Bodies.\n\n Herbert Dingle (1980) The Twins Paradox of Relativity.\n\n Hafele, Joseph C.; Keating, Richard E. (July 14, 1972). \"Around-the-World Atomic Clocks: Predicted Relativistic Time Gains\". Science 177 (4044): 166–168. Bibcode 1972Sci...177..166H. doi:10.1126/science.177.4044.166. PMID 17779917.\n\n Hafele, Joseph C.; Keating, Richard E. (July 14, 1972). \"Around-the-World Atomic Clocks: Observed Relativistic Time Gains\". Science 177 (4044): 168–170. Bibcode 1972Sci...177..168H. doi:10.1126/science.177.4044.168. PMID 17779918.\n\n Murray, W.A. Scott, \"If you want to know the time...\" Wireless World. 92 No.16 28-31 1986.\n\n Geoffrey Builder (1957) Ether and Relativity.\n\n Hendrick Lorentz (1904) Electromagnetic phenomena in a system moving with any velocity smaller than that of light.\n\nAppendix A\n\nPost’s Kinematic explanation of the Sagnac Effect:\n\nE.J. Post uses the stationary observer position where the speed of light is assumed to be c, and then works backwards to try to\n\nshow that the moving observer’s measured result is the same. From the point of view of the stationary observer:\n\nPosition C = origin/detector\n\nDs’ = distance C has rotated clockwise by the time the counter-clockwise beam reaches it at C’\n\nDs” = distance C has rotated by the time the clockwise beam reaches it at C”\n\nc = free space speed of light.", null, "Time difference for each path t’ = counter-clockwise, t” = clockwise", null, "From the stationary observer’s point of view, the time difference is then:", null, "The time interval of the rotating beam splitter is then computed to be:", null, "Where g corresponds to the time dilation factor due to the rotational motion.  However, does the relativistic correction to the Sagnac equation actually make a difference in the actual result? When we consider the rotating Sagnac interferometer as a device experiencing time dilation, we need to consider first which parts are clocks and what the actual outcome time dilation will have on the measuring system, the detector. Let’s look at the Sagnac interferometer as three parts:\n\na)      A source such as a laser, which is a light clock subject to velocity time dilation\n\nb)     An optical path composed of wavelengths of the emitted light\n\nc)      A detector, which measures the phase difference between the wave-fronts of the two returning beams.\n\nThe source (for example a laser) is a true clock – lets assume that this clock experiences a time dilation proportional to\n\nt’ = to/SQRT(1-WR^2/c^2)  = to * b .  then,\n\nDt = to * b - to   =  (2L- 2L/b)/C                                 (1)\n\nThis will inevitably mean that the rotating source must emit light at a reduced frequency, proportional to f’ = f * SQRT(1-WR^2/c^2). However, since C = f*l, the wavelengths in the optical path must expand due to the lower frequency. This means that there will be fewer wavelengths passing the detector per unit time.  The difference will be in the amount of time a given wave-front is now behind compared to the unaffected laser situation:\n\nDt = (2L- 2L/b)/C                                                    (2)\n\nSo the wave-fronts are behind by the amount Dt. But as we showed in 1, the time is also behind by the same amount Dt since the duration of each second is larger due to the clock frequency change. The result is that the two effects cancel out, and the same wave-front arrives at the detector at the same elapsed time regardless of the time dilation, from the standpoint of the observer on the disk. The relativistic correction would seem thereby to be unnecessary. As Post says, the result is too small to be experimentally distinguishable, and the common practice is to use a g factor = 1.\n\nSo this in and of itself does not contradict the notion that the measured speed of light by an observer on the disk is C +/- WR.\n\nPost says himself that “the circulation times t1’ and t2’ in the rotating frame are obtained by integrating f’ from 0-2p and 0- -2p, respectively:", null, "And then it is easy to see that in the rotating frame:\n\nDt = L/(C+WR) – L/(C- WR) = 4AW/c^2                         (3)\n\nPost quotes the method of Langevin from his 1921 and 1937 papers. In his 1921 paper, Langevin argued that because the Sagnac effect was first order, there could be no distinction between Newtonian or relativistic derivations and therefore the effect was in accordance with relativity theory. Dufour and Prunier took Langevin to task on this. Dufour and Prunier’s work is discussed by Kelly in his book Challenging Modern Physics: Questioning Einstein's Relativity Theories”.    In their experiments they created Sagnac interferometers that were composites of moving and stationary paths, including stationary sources and stationary detectors. This was essentially to test if the relativistic approach could be distinguished from the classical approach. In particular, they focused on the mathematical treatment of Langevin, whereby the position placement of the observer on the disk is arbitrary. Using Langevin’s calculation:\n\nDt = L/C + 2WA/C^2                      (4)\n\nThey predicted that in a Sagnac interferometer where only part of the light path was rotating and the other part stationary, the classical equation would call for a different area A than the relativistic equation for calculating the result.  In all cases of this experimental test, the Sagnac effect was the same. This overturned Langevin’s analysis, and in 1937, he had to revise his explanation, as pointed out by Kelly:\n\n“In his final essay on the subject in 1937, Langevin proposed that the results published that year by Dufour and Prunier showed that one had to assume either (a) the light speed varied to c + wr in one direction and c – wr in the other direction, or (b) the time aboard the spinning apparatus had to change by a factor of +/-2wA/c2 in either direction. Indeed, Langevin went as far as to say that assuming (a), “we find, by a very simple and very general reasoning, the formula for the difference of the times of the path of the two light beams in the Sagnac experiment.” He also stated that “the paths are of unequal length because of the inequality of the speed of prorogation”.\n\nThe proposition (b) though is untenable because if this were true then when the light beam passed back to the moving detector, the local time from each direction would be out of synchronization, meaning that the clocks cannot be counting real time and that the effective time dilation is meaningless. This was also pointed out by Herbert Ives in his 1938 paper criticizing Langevin. Ives says about the absurdity of Langevin’s proposition (b):\n\n” There are of course not merely two clocks, but an infinity of clocks, where we include those that could be transported at finite speeds, and around other paths. As emphasized previously, the idea of “local time” is untenable, what we have are clock readings. Any number of clock readings at the same place are physically possible, depending on the behaviour and history of the  clocks used. More than one “time” at one place is a physical absurdity. “\n\nThe only explanation left, is Langevin’s proposition a) that the light speed varies by C+/-wr in one or the other direction around the disk, consistent with Dufour and Prunier’s experimental results.\n\nReferences:\n\n Paul Langevin – (1937) On the Experiment of Sagnac. Compt. Rend. 205, 51.\n\n[15} Paul Langevin – (1921) Compt. Rend. 173, 831.\n\n A. Dufour, F. Prunier,\n\n E. J. Post, Sagnac Effect, Review of Modern Physics, 39,2," ]
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https://codescracker.com/php/php-break-continue-keyword.htm
[ "Tutorials\nExamples\nTests\n\n# break and continue Statement in PHP\n\n1. break\n2. continue\n\n## PHP break Statement\n\nThe PHP break keyword or statement, is used when we need to exit from the current loop or switch structure. For example:\n\n```<?php\nfor(\\$i=1; \\$i<10; \\$i++)\n{\nif(\\$i>5)\nbreak;\n\necho \\$i;\necho \"<BR>\";\n}\n?>```\n\nThe snapshot given below shows the output produced by this PHP example on break statement:\n\nIn above example, when the value of \\$i becomes greater than 5, that is when the value of \\$i becomes 6, then the condition \\$i>5 or 6>5 evaluates to be true, therefore, the program flow goes inside the if block and the break; statement gets executed, that skips the remaining execution of the current for loop. Therefore, we only get the output from 1 to 5.\n\nThe break statement skips the execution of remaining statement(s) (if any) along with all the remaining iterations of the current loop.\n\nThe break statement can also be used to stop the execution of an infinite loop (the loop whose condition always evaluates to be true). For example:\n\n```<?php\n\\$num = 10;\nwhile(true)\n{\necho \\$num, \"<BR>\";\n\nif(\\$num==12)\nbreak;\n\n\\$num++;\necho \\$num, \"<BR>\";\n}\n?>```\n\nThe output of this PHP example should be:\n\nSee the condition of the loop, it is true, that obviously evaluates to be true for ever. Therefore, using the break statement, the execution of this loop ends.\n\nThe dry run of above example goes like:\n\n• The value 10 initialized to the variable \\$num using the first statement\n• Now the execution of the while loop started\n• Since the condition of the loop is true, therefore, the loop goes for its infinite execution\n• At first iteration, using the echo, the value of \\$num gets printed along with a line break using BR tag\n• Since the condition \\$num==12 or 10==12 evaluates to be false, therefore program flow will not go inside the if block to execute the break statement\n• The value of \\$num gets incremented by 1. Now \\$num=11\n• The value of \\$num, again gets printed along with a line break\n• Now the second iteration of the loop starts\n• This process continues, until the value of \\$num becomes 12\n• When the value of \\$num becomes 12, then the condition \\$num==12 (of if) evaluates to be true, and the break; statement gets executed, that forces the while loop to stop its execution\n\n## PHP continue Statement\n\nThe PHP continue statement or keyword is used, when we need to skip the execution of remaining statement(s) for the current iteration of the current loop, and continue for the next iteration of the current loop. For example:\n\n```<?php\nfor(\\$i=1; \\$i<10; \\$i++)\n{\nif(\\$i==5)\ncontinue;\n\necho \\$i;\necho \"<BR>\";\n}\n?>```\n\nThe output of above PHP example on continue statement, is:\n\nIn above example, when the value of \\$i becomes 5, then the condition \\$i==5 evaluates to be true. Therefore the continue; statement gets executed, that forces to jump and continue the execution of next iteration of the current loop. That is, the two echo statement gets skipped for that iteration.\n\n## Difference Between break and continue in PHP\n\nThe break is used to skip all the remaining execution and iteration of the current loop. Whereas the continue is used to skip the execution of the remaining statement(s) of the current iteration of the current loop. For example:\n\n```<?php\necho \"<h2>Using the break Statement</h2>\";\nfor(\\$i=1; \\$i<=5; \\$i++)\n{\nif(\\$i==3)\nbreak;\necho \\$i, \"<BR>\";\n}\n\necho \"<h2>Using the continue Statement</h2>\";\nfor(\\$i=1; \\$i<=5; \\$i++)\n{\nif(\\$i==3)\ncontinue;\necho \\$i, \"<BR>\";\n}\n?>```\n\nThe output of above PHP example, is:\n\nPHP Online Test\n\n« Previous Tutorial Next Tutorial »" ]
[ null ]
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https://cs50.stackexchange.com/questions/33345/pset4-speller-malloc-and-calloc-usage
[ "# PSET4 (Speller) malloc and calloc usage\n\nAfter getting my code to work, i looked if there were any memory leaks using valgrind. It found no memory leaks, however it returned these two errors:\n\n``````==7109== Conditional jump or move depends on uninitialised value(s)\n==7109== at 0x4C32D08: strlen (in /usr/lib/valgrind/vgpreload_memcheck-amd64-linux.so)\n==7109== by 0x40138D: check (dictionary.c:154)\n==7109== by 0x400C39: main (speller.c:112)\n==7109== Uninitialised value was created by a heap allocation\n==7109== at 0x4C2FB0F: malloc (in /usr/lib/valgrind/vgpreload_memcheck-amd64-linux.so)\n==7109== by 0x400914: main (speller.c:40)\n``````\n\nAfter searching online, I discovered the calloc method that initializes the pointer given by malloc to 0. After using it in my code, I got no more errors in valgrind, but i don't understand why? I don't understand what `calloc()` is doing that `malloc()` isn't. I understand that calloc does the same as malloc but initializes the memory block to 0.\n\n``````// Loads dictionary into memory, returning true if successful else false\n{\n// Initialize hash table\nfor (int i = 0; i < N; i++)\n{\nhashtable[i] = NULL;\n}\n\n// Open dictionary\nFILE *file = fopen(dictionary, \"r\");\nif (file == NULL)\n{\nreturn false;\n}\n\n// Buffer for a word\nchar word[LENGTH + 1];\n\n// Insert words into hash table\nwhile (fscanf(file, \"%s\", word) != EOF)\n{\n// Create a node with the current word\nnode *element = malloc(sizeof(node)); // allocate memory for a new node\n\n// Check if node was created sucessfully\nif (element == NULL)\n{\nprintf(\"Error creating node\\n\");\nreturn false;\n}\n// Initialize current node\nelement->next = NULL;\n``````\n\nnode structure:\n\n``````typedef struct node\n{\nchar word[LENGTH + 1];\nstruct node *next;\n}\nnode;\n``````\n\n``````// Adds a word to a node's value (word)\nint add_word(node *element, const char *word)\n{\nif (element == NULL)\n{\nprintf(\"Node is invalid\");\nreturn 1;\n}\n\nfor (int i = 0; word[i] != '\\0'; ++i)\n{\nelement->word[i] = word[i];\n}\n\n// Success\nreturn 0;\n}\n``````\n• What does your node structure look like? What does `add_word` do? Does it use the `next` member before you set it to `NULL`? Aug 30 '19 at 13:29\n• Sorry forgot to include those. I've edited my original quesiton Aug 30 '19 at 14:30\n• `next` obviously was the wrong thought by me, as `strlen` was involved (totally missed that). Should have guessed \"missing null terminator\" immediately. Thanks for adding the code, made it easier to see the mistake. Aug 30 '19 at 14:58\n\nUsing `calloc` would explicitly zero all the bytes, including the one later serving as null terminator, so this mistake would no longer cause `valgrind` to complain, it's just writing more zeroes than required.\nBTW, there's a function that essentially does what your `add_word` does, and it writes a null terminator: `strcpy`." ]
[ null ]
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https://www.numbersaplenty.com/5072059
[ "Search a number\nBaseRepresentation\nbin10011010110010010111011\n3100112200120001\n4103112102323\n52244301214\n6300413431\n761053226\noct23262273\n910480501\n105072059\n112954793\n121847277\n131087825\n149605bd\n156a2c74\nhex4d64bb\n\n5072059 has 4 divisors (see below), whose sum is σ = 5077800. Its totient is φ = 5066320.\n\nThe previous prime is 5072057. The next prime is 5072063. The reversal of 5072059 is 9502705.\n\nIt is a semiprime because it is the product of two primes, and also a brilliant number, because the two primes have the same length, and also an emirpimes, since its reverse is a distinct semiprime: 9502705 = 51900541.\n\nIt is a cyclic number.\n\nIt is not a de Polignac number, because 5072059 - 21 = 5072057 is a prime.\n\nIt is a Duffinian number.\n\nIt is not an unprimeable number, because it can be changed into a prime (5072057) by changing a digit.\n\nIt is a Kaprekar number, because its square (25725782499481) can be partitioned into two parts whose sum is 5072059.\n\nIt is a pernicious number, because its binary representation contains a prime number (13) of ones.\n\nIt is a polite number, since it can be written in 3 ways as a sum of consecutive naturals, for example, 1234 + ... + 3415.\n\nIt is an arithmetic number, because the mean of its divisors is an integer number (1269450).\n\nAlmost surely, 25072059 is an apocalyptic number.\n\n5072059 is a deficient number, since it is larger than the sum of its proper divisors (5741).\n\n5072059 is a wasteful number, since it uses less digits than its factorization.\n\n5072059 is an odious number, because the sum of its binary digits is odd.\n\nThe sum of its prime factors is 5740.\n\nThe product of its (nonzero) digits is 3150, while the sum is 28.\n\nThe square root of 5072059 is about 2252.1232204300. The cubic root of 5072059 is about 171.8151407639.\n\nThe spelling of 5072059 in words is \"five million, seventy-two thousand, fifty-nine\".\n\nDivisors: 1 1091 4649 5072059" ]
[ null ]
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https://oeis.org/A127885
[ "The OEIS Foundation is supported by donations from users of the OEIS and by a grant from the Simons Foundation.", null, "Hints (Greetings from The On-Line Encyclopedia of Integer Sequences!)\n A127885 a(n) = minimal number of steps to get from n to 1, where a step is x -> 3x+1 if x is odd, or x -> either x/2 or 3x+1 if x is even; or -1 if 1 is never reached. 14\n 0, 1, 7, 2, 5, 8, 16, 3, 11, 6, 14, 9, 9, 17, 17, 4, 12, 12, 20, 7, 7, 15, 15, 10, 23, 10, 23, 10, 18, 18, 31, 5, 18, 13, 13, 13, 13, 21, 26, 8, 21, 8, 21, 16, 16, 16, 29, 11, 16, 16, 24, 11, 11, 24, 24, 11, 24, 19, 24, 19, 19, 32, 32, 6, 19, 19, 27, 14, 14, 14, 27, 14, 27, 14, 14, 22, 22, 27, 27, 9, 22, 22, 22, 9, 9, 22, 22, 17, 22, 17, 30, 17, 17, 30, 30, 12, 30, 17, 17, 17 (list; graph; refs; listen; history; text; internal format)\n OFFSET 1,3 COMMENTS In contrast to the \"3x+1\" problem (see A006577), here you are free to choose either step if x is even. See A125731 for the number of steps in the reverse direction, from 1 to n. REFERENCES M. J. Halm, Sequences (Re)discovered, Mpossibilities 81 (Aug. 2002), p. 1. LINKS David Applegate, Table of n, a(n) for n = 1..1000 FORMULA a(1) = 0; and for n > 1, if n is odd, a(n) = 1 + a(3n+1), and if n is even, a(n) = 1 + min(a(3n+1),a(n/2)). [But with a similar caveat as in A257265] - Antti Karttunen, Aug 20 2017 EXAMPLE Several early values use the path: 6 -> 3 -> 10 -> 5 -> 16 -> 8 -> 4 -> 2 -> 1. The first path where choosing 3x+1 for even x helps is: 9 -> 28 -> 85 -> 256 -> 128 -> 64 -> 32 -> 16 -> 8 -> 4 -> 2 -> 1. MAPLE # Code from David Applegate: Be careful - the function takes an iteration limit and returns the limit if it wasn't able to determine the answer (that is, if A127885(n, lim) == lim, all you know is that the value is >= lim). To use it, do manual iteration on the limit. A127885 := proc(n, lim) local d, d2; options remember; if (n = 1) then return 0; end if; if (lim <= 0) then return 0; end if; if (n > 2 ^ lim) then return lim; end if; if (n mod 2 = 0) then d := A127885(n/2, lim-1); d2 := A127885(3*n+1, d); if (d2 < d) then d := d2; end if; else d := A127885(3*n+1, lim-1); end if; return 1+d; end proc; MATHEMATICA Table[-1 + Length@ NestWhileList[Flatten[# /. {k_ /; OddQ@ k :> 3 k + 1, k_ /; EvenQ@ k :> {k/2, 3 k + 1}}] &, {n}, FreeQ[#, 1] &], {n, 126}] (* Michael De Vlieger, Aug 20 2017 *) PROG (PARI) { A127885(n) = my(S, k); S=[n]; k=0; while( S!=1, k++; S=vecsort( concat(apply(x->3*x+1, S), apply(x->x\\2, select(x->x%2==0, S) )), , 8);  ); k } /* Max Alekseyev, Sep 03 2015 */ CROSSREFS A127886 gives the difference between A006577 and this sequence. Cf. A006577, A125731, A127887, A125195, A125686, A125719, A261870. Cf. A290100 (size of the final set when using Alekseyev's algorithm). Cf. also A257265. Sequence in context: A072761 A337357 A340420 * A006577 A337150 A280234 Adjacent sequences:  A127882 A127883 A127884 * A127886 A127887 A127888 KEYWORD nonn AUTHOR David Applegate and N. J. A. Sloane, Feb 04 2007 EXTENSIONS Escape clause added to definition by N. J. A. Sloane, Aug 20 2017 STATUS approved\n\nLookup | Welcome | Wiki | Register | Music | Plot 2 | Demos | Index | Browse | More | WebCam\nContribute new seq. or comment | Format | Style Sheet | Transforms | Superseeker | Recent\nThe OEIS Community | Maintained by The OEIS Foundation Inc.\n\nLast modified June 16 14:16 EDT 2021. Contains 345057 sequences. (Running on oeis4.)" ]
[ null, "https://oeis.org/banner2021.jpg", null ]
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https://discourse.julialang.org/t/implementing-gibbs-sampler-for-ode-parameter-estimation/99132
[ "# Implementing Gibbs sampler for ODE parameter estimation\n\nI would like to use a Gibbs sampler for a model that looks like this (in Turing.jl):\n\n``````@model function fit_burger2(data::AbstractVector, prob)\n# Prior distributions\nσᵣ = 1.5e6\nAₘ ~ truncated(Normal(0.1, 0.05); lower=0.01, upper=1.1)\nQₘ ~ Uniform(1e3,1e6)\nn ~ Uniform(0.7,4.5)\npₘ ~ truncated(Normal(0.36, 0.2); lower=0.3, upper=0.7)\nEₖ ~ Uniform(1e9,1e10)\nEₘ ~ Uniform(1e8,9e9)\nη ~ truncated(Normal(1e11, 0.5e9); lower=1e9, upper=1e12)\ndₛₛ ~ truncated(Normal(2.4e-4, 1e-4); lower=1e-5, upper=1e-3)\nFₛₛ ~ truncated(Normal(0.5, 0.1); lower=0.3, upper=0.7)\nε₁ ~ Uniform(0.01,1.0)\nε₂ ~ truncated(Normal(0.16, 0.07); lower=0.01, upper=0.9)\nS = 1e-4\nR = 8.3145\nT = 263\n\n# Simulate the model\np = [Aₘ, Qₘ, n, pₘ, R, T, Eₖ, Eₘ, η, dₛₛ, Fₛₛ, ε₁, ε₂, S]\npredicted = solve(prob, Rosenbrock23(); p=p, saveat=df2.time, save_idxs=3)\n\n# Observations\ndata ~ MvNormal(predicted.u, σᵣ^2 * I)\n\nreturn nothing\nend\n``````\n\nI wonder what would be the process to do it, currently, I’m using regular Metropolis-Hastings:\n\n``````chain2 = sample(model3, MH(), MCMCSerial(), 120_000, 2;\nprogress=true,\nI was trying to find examples/tutorials that use Gibbs for estimating ODE parameters but couldn’t find any. So any hint to throw me in the right direction would be highly appreciated", null, "What sort of Gibbs sampling are you thinking?", null, "" ]
[ null, "https://emoji.discourse-cdn.com/twitter/slight_smile.png", null, "https://emoji.discourse-cdn.com/twitter/slight_smile.png", null ]
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https://www.armoredpenguin.com/wordsearch/Data/2019.05/2804/28041826.152.html
[ "", null, "", null, "", null, "", null, "", null, "", null, "", null, "", null, "### Talbot 2R1\n\nPersuasive techniques\n\n k x d o t c i e r o s p a a t e n e n s p e e c h m a r k s e a m i n p g o s c i t s i t a t s e n o n n a s s i s s g n p c t i i g o t e e i e s i t a n e x p e r t w i t n e s s h e e s s s x m i g s i t n u s r r o l v r e e i a e y n l a p s e l h r g r i i u v t g t i h r a e t i r t q e p t p q i t g l s e r l v s i s f u p s a r l t e e a g f i l t p e o o o c s r o a a t r g a m s i a n r o n t n n a n c l o a c i i n t t n t a r a s r p o i r x t s i g t e e h r n e t t o m u r e s i o e e n r r e p e t i t i o n o p t o n r c y a t s i t t o t t c s t u e n n e m o t i v e l a n g u a g e s e i s i s p i l l e i p f i s n s h c n n n p n n o i t a m a l c x e s r g i a t t o l n s r o h p a t e m e u\nfacts\ninternalrhyme\nexaggeration\nalliteration\nstatistics\npronouns\nrepetition\nexaggeration\nrhetoricalquestion\ncomparatives\nellipsis\nsuperlatives\nexclamation\nexpertwitness\nquotation\nshortsentences\nlisting\npatterning\ndescription\npatternofthree\nspeechmarks\nemotivelanguage\nmetaphors\nsimiles\n\nSome of the puzzles that people list for the public get indexed by the search engines (like Google). Some people find those puzzles and cannot figure out how to make a puzzle of their own. So this page now has the navigation sidebar.\n\nThere are now buttons on the puzzle so that you can get a clean page, in either HTML or PDF, that you can use your browser's print button to print. The PDF format allows the web site to know how large a printer page is, and the fonts are scaled to fill the page. The PDF takes awhile to generate. Don't panic!", null, "Web armoredpenguin.com" ]
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https://socratic.org/questions/what-is-the-polar-form-of-17-15
[ "# What is the polar form of ( 17,15 )?\n\n$\\left(\\sqrt{514} , {41.423665625}^{\\circ}\\right)$\n\n#### Explanation:\n\nGiven $x = 17$ and $y = 15$\n\n$r = \\sqrt{{x}^{2} + {y}^{2}} = \\sqrt{{17}^{2} + {15}^{2}} = \\sqrt{289 + 225} = \\sqrt{514}$\n\n$\\theta = {\\tan}^{-} 1 \\left(\\frac{y}{x}\\right) = {\\tan}^{-} 1 \\left(\\frac{15}{17}\\right) = {41.423665625}^{\\circ}$\n\nGod bless....I hope the explanation is useful." ]
[ null ]
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https://www.findfilo.com/maths-question-answers/evaluate-lim-xrarr2-x-2-log-x-2-where-represents-trhb
[ "", null, "Evaluate : (\"lim\")_(xrarr2^+) ([x-2])/(\"log\"(x-2)) , where [.] | Filo", null, "", null, "Class 11\n\nMath\n\nCalculus\n\nLimits And Derivatives", null, "558\n\nEvaluate : , where [.] represents the greatest integer function.", null, "558", null, "Connecting you to a tutor in 60 seconds." ]
[ null, "https://www.facebook.com/tr", null, "https://www.findfilo.com/images/logo.svg", null, "https://www.findfilo.com/images/icons/navbar.png", null, "https://www.findfilo.com/images/icons/view.svg", null, "https://www.findfilo.com/images/icons/view.svg", null, "https://www.findfilo.com/images/common/mobile-widget.png", null ]
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https://www.xsanisty.com/project/calx/comment-page-1/
[ "## What is jQuery-Calx?\n\njQuery Calx is powerful yet easy to use jQuery plugin for building a calculation form or calculation table, it’s parse provided formula and do calculation based on it, scan the form change and update the result automatically, format plain number into currency format, ordinal number, etc.\n\njQuery Calx was designed to allow user to easily configure their calculation form or calculation table, you may define as simple as (\\$A+\\$B) formula to the complex one such as PMT formula `(\\$I*\\$P*((1 + \\$I)^\\$N)) / (1 - ((1 + \\$I)^\\$N)).`\n\n## What operators are supported by jQuery-Calx\n\nCurrently, jQuery-Calx support standard calculation operator, nested parentheses, aggregation, etc.\n\n• [+,-,/,*] : standard addition, subtraction, divide and multiply. sample: (\\$A+\\$B/\\$C), etc\n• [^, SQRT] : power and square root also possible for calculation. sample: (\\$A^\\$B), (SQRT(\\$C)), etc\n• [MOD]: modulus operator will find the remainder of division of one number by another. sample(\\$A MOD \\$B);\n• [MAX,MIN,SUM,AVG] : simple aggregation function, it’s only work when the element id mimic the excel cell index such as A1, C4, etc. sample: (MAX(\\$A1,\\$B6)).\n• [<,<=,>=,=,<>]: comparative operator, only usable with IF statement. sample: IF(\\$A < \\$B, \\$A, \\$B).\n• [AND, OR]: logic operator, used within IF statement. IF(\\$A < \\$B AND \\$B <> 0 AND \\$A <> 0, \\$A, \\$B).\n• [IF]: simple logic condition, you may nest the condition as deep as you want IF([condition],[true],[false]). sample: IF(\\$A <> \\$B, \\$C, \\$D), IF(\\$A = 1, \\$B, IF(\\$A = 2, \\$B, \\$C))\n\nfor detailed documentation about formula operator and how to use it, please refer to ‘Using Formula with Calx‘ page.\n\n## What is it useful for?\n\njQuery-Calx is useful when you need to build web based calculator, item and discount calculator, calculation table, etc. It also useful when you need to convert excel spreadsheet calculator into web based calculator.\n\n## Contribute?\n\nIf you found a bug, mis-calculation, strange behavior, as well as expect new feature, please kindly report an issue on github issue tracker or contact me directly here.\n\nIf you found a bug and able to fix it yourself, or add new feature to the calx, please kindly fork my  github repo and send a pull request.\n\nIf you found this plugin is useful for you but can not contribute any code or reporting any issues, you can also support the development by treat us some lunch 🙂", null, "" ]
[ null, "https://www.paypalobjects.com/en_US/i/scr/pixel.gif", null ]
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https://www.geeksforgeeks.org/rotate-instructions-in-8085/
[ "# ROTATE Instructions in 8085\n\nROTATE is a logical operation of 8085 microprocessor. It is a 1 byte instruction. This instruction does not require any operand after the opcode. It operates the content of accumulator and the result is also stored in the accumulator. The Rotate instruction is used to rotating the bits of accumulator.\n\nTypes of ROTATE Instruction:\nThere are 4 categories of the ROTATE instruction: Rotate accumulator left (RLC), Rotate accumulator left through carry (RAL), Rotate accumulator right (RRC), Rotate accumulator right through carry (RAR). Among these four instructions; two are for rotating left and two are for rotating right. All of them are explain briefly in the following sections:\n\n1. Rotate accumulator left (RLC) –\nIn this instruction, each bit is shifted to the adjacent left position. Bit D7 becomes D0. Carry flag CY is modified according to the bit D7. For example:-\n\n```A = D7 D6 D5 D4 D3 D2 D2 D0\n\n//before the instruction\nA = 10101010; CY=0\n\n//after 1st RLC\nA = 01010101; CY=1\n\n//after 2nd RLC\nA = 10101010; CY=0 ```\n2. Rotate accumulator left through carry (RAL) –\nIn this instruction, each bit is shifted to the adjacent left position. Bit D7 becomes the carry bit and the carry bit is shifted into D0. Carry flag CY is modified according to the bit D7. For example:\n\n```A = D7 D6 D5 D4 D3 D2 D2 D0\n\n//before the instruction\nA = 10101010; CY=0\n\n//after 1st RAL\nA = 01010100; CY=1\n\n//after 2nd RAL\nA = 10101001; CY=0 ```\n3. Rotate accumulator right (RRC) –\nIn this instruction, each bit is shifted to the adjacent right position. Bit D7 becomes D0. Carry flag CY is modified according to the bit D0. For example:\n\n```A = D7 D6 D5 D4 D3 D2 D2 D0\n\n//before the instruction\nA = 10000001; CY=0\n\n//after 1st RRC\nA = 11000000; CY=1\n\n//after 2nd RRC\nA = 01100000; CY=0 ```\n4. Rotate accumulator right through carry (RAR) –\nIn this instruction, each bit is shifted to the adjacent right position. Bit D0 becomes the carry bit and the carry bit is shifted into D7. Carry flag CY is modified according to the bit D0. For example:\n\n```A = D7 D6 D5 D4 D3 D2 D2 D0\n\n//before the instruction\nA = 10000001; CY=0\n\n//after 1st RAR\nA = 01000000; CY=1\n\n//after 2nd RAR\nA = 10100000; CY=0 ```\n\nApplications of ROTATE Instructions:\nThe ROTATE instructions are primarily used in arithmetic multiply and divide operations and for serial data transfer. For example:\n\n```If A is 0000 1000 = 08H\n\n1. By rotating 08H right : A = 0000 0100 = 04H\nThis is equivalent to dividing by 2.\n\n2. By rotating 08H left : A = 0001 0000 = 10H\nThis is equivalent to multiplying by 2. ```\n\nHowever, these procedures are invalid when logic 1 is rotated left from D7 to D0 or vice versa. For example, if 80H is rotated left it becomes 01H.", null, "My Personal Notes arrow_drop_up", null, "Check out this Author's contributed articles.\n\nIf you like GeeksforGeeks and would like to contribute, you can also write an article using contribute.geeksforgeeks.org or mail your article to [email protected]. See your article appearing on the GeeksforGeeks main page and help other Geeks.\n\nPlease Improve this article if you find anything incorrect by clicking on the \"Improve Article\" button below.\n\nImproved By : SakshamGoyal" ]
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https://www.growingiq.com/dlt-collection/Ms.-Joan-/level2/saturday/307a6ba9-974b-4dc1-bde2-cf645dff6580
[ "## Ms. Joan\n\n### Target 1​\n\n###### Lesson Type:\n\nNew\n\nNumber Operation\n\n:\n\nComputation\n\nSubtract within 50.\n\n###### 1:\n\nRecognize the arithmetic symbols in a subtraction equation: minus (-) and equals (=).\n\n###### 2:\n\nUnderstand that subtraction means to take a quantity away from another quantity.\n\n###### 3:\n\nUnderstand that subtracting numbers means making a smaller number.\n\nKindergarden\n\n###### Vocabulary:\n\ntake away, plus, minus, equals\n\nActivities:\n\nStudents subtract using counters\n\n• Students subtracted using drawing then wrote a number sentence to represent it.\n• To get the difference of a given equation. students counted backwards. To help students visualize counting backwards, students were asked to write the numbers.", null, "", null, "", null, "", null, "### Home Exploration\n\n###### Guiding Questions:", null, "## Absent Students:\n\n### Target 2\n\n:\n\n###### 1:\n\nIdentify the correct arithmetic process based on the information presented in word problems.\n\n###### 2:\n\nFind and use the needed information in a word problem in order to solve.\n\nKindergarden\n\n###### Vocabulary:\n\naddition, subtraction, take-away, put together, altogether\n\nActivities:\n\n• Word problems were read to the students. Students identified if the word problem is addition or subtraction.\n• They were also given different ways to show how to solve the word problem. (drawing, number line, tens frame, number sentence)", null, "", null, "", null, "", null, "### Home Exploration\n\n###### Guiding Questions:", null, "### Target 3\n\n:\n\n###### Vocabulary:\n\nActivities:", null, "", null, "", null, "", null, "### Home Exploration", null, "" ]
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https://www.examrace.com/ACET/ACET-Free-Study-Material/Statistics/Statistics-Index-Numbers.html
[ "ACET: Index Numbers\n\nGlide to success with Doorsteptutor material for ACET : get questions, notes, tests, video lectures and more- for all subjects of ACET.\n\nhow do you define index numbers? Narrate the nature and types of index number with adequate example.\n\naccording to Croxton and Cowden index numbers are devices for measuring difference sin the magnitude of a group of related\n\nAccording to Morris Hamburg in its simplest form an index number is nothing more than a relative which express the relationship between two figures, where one figure is used as a base. According to M L Berenson and D M Levine generally speaking, index number measure the size or magnitude of some object at particular point in time as a percentage of some base or reference object in the past. According to Richard. I. Levin and David S. Rubin an index number is a measure how much a variable changes over time\n\nNature of Index Number\n\n1. Index numbers are specified average used for comparison in situation where two or more series are expressed in different units or represent different items. E. g. Consumer price index representing prices of various items or the index of industrial production representing various commodities produced.\n2. Index number measure the net change in a group of related variable over a period of time.\n3. Index number measure the effect of change over a period of time, across the range of industries, geographical regions or countries.\n4. The consumption of the index number is carefully planned according to the purpose of their computation, collection of data and application of appropriate method, assigning of correct weightages and formula.\n\nTypes of Index Numbers\n\nPrice index numbers: A price index is any single number calculated from an array of prices and quantities over a period. Since not all prices and quantities of purchases can be recorded, a representative sample is used instead. Price are generally represented by p in formulae. These are also expressed as price relative, defined as follows\n\nPrice relative = (current years price/base years price) ⚹ 100, i.e.. . = (p1/p0) ⚹ 100 any increses in price index amounts to corresponding decreses in purchasing power of the rupees or other affected currency. Quantity index number a quantity index number measures how much the number or quantity of a variable changes over time. Quantities are generally represented as q in formulae. Value index number: a value index number measures changes in total monetary worth, that is, it measure changes in the rupee value of a variable. It combines price and quantity changes to present a more informative index. Composite index number: a single index number may reflect a composite, or group, of changing variable. For instance, the consumer price index measures the general price level for specific goods and service in the economy. These are also known as index numbers. In such cases the price-relative with respect to a selected base are determined separately for each and their statistical average is computed\n\nwhat are the importance of index numbers used in Indian economy. Explain index numbers of industrial production\n\nImportance of Index Numbers Used in Indian Economy\n\nCost of living index or consumer price index\n\nCost of living index number or consumer price index, expressed as percentage, measure the relative amount of money necessary to derive equal satisfaction during two periods of time, after taking into consideration the fluctuations of the retail prices of consumer goods during these periods. This index is relevant to that real wages of workers are defined as (actual wages/cost of living index) ⚹ 100. Generally the list of items consumed varies for different classes of people (rich, middle, class, or the poor) at the same place of residence. Also people of the same class belonging to different geographical regions have different consumer habits. Thus the cost of living index always relates to specific class of people and a specific geographical area, and it help in determining the effect of changes in price on different classes of consumers living in different areas. The process of construction of cost of living index number is as follows\n\n1. Obtain decision about class of people for whom the index number is to be computed, for instance, the industrial personnel, officers or teachers etc. Also decide on the geographical area to be covered.\n2. Conduct a family budget inquiry covering the class of people for whom the index number is to be computed. The enquiry should be conducted for the base year by the process of random sampling. This would give information regarding the nature, quality and quantities of commodities consumed by an average family of the class and also the amount spent on different items of consumption.\n3. The item on which the information regarding money spent is to be collected are food (rice, wheat, sugar, milk, tea etc) , clothing, fuel and lighting, housing and miscellaneous items.\n4. Collect retail prices in respect of the items from the localities in which the class of people concerned reside, or from the markets where they usually make their purchases.\n5. as the relative importance of various items for different classes of people is not the same, the price or price relative are always weighted and therefore, the cost of living index is always a weighted index.\n6. The percentage expenditure on an item constitutes the weight of the item and the percentage expenditure in the five groups constitutes the group weight.\n7. Separate index number are first of all determined for each of the five major groups, by calculating the weighted average of price-relatives of the selected items in the group.\n\nIndex Number of Industrial Production\n\nThe index number of industrial production is designed to measure increase or decrease in the level of industrial production in a given period of time compared to some base periods. Such an index measures changes in the quantities of production and not their values. Data about the level of industrial output in the base period and in the given period is to be collected first under the following heads\n\n• Textile industries to include cotton, woolen, silk etc.\n• Mining industries like iron ore, iron, coal, copper, petroleum etc.\n• Metallurgical industries like automobiles, locomotive, aero planes etc\n• Industries subject to excise duties like sugar, tobacco, match etc.\n• Miscellaneous like glass, detergents, chemical, cement etc.\n\nExplanation\n\nThe figure of output for a various industries classifies above are obtained on a monthly, quarterly or yearly basis. Weights are assigned to various industries on the basis of some criteria such as capital invested turnover, net output, production etc. Usually the weights in the index are based on the values of net output of different industries. The index of industrial production is obtained by taking the simple mean or geometric mean of relatives. When the simple arithmetic mean is used the formula for constructing the index is as follows.\n\nIndex of industrial production = (100/w) ⚹ (q1/q0) ⚹ w = (100/w) ⚹ I ⚹ w\n\n• Where q1 = quantity produced in a given period\n• Q0 = quantity produced in the base period\n• W = relative importance of different outputs\n• I = (q1/q0) = index for respective commodity\n\nDeveloped by:" ]
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