content
stringlengths 0
1.88M
| url
stringlengths 0
5.28k
|
---|---|
For this assignment, please create a Requirements Management Plan (RMP), and read Project Scenario carefully.
The RMP is a key document in managing the scope of your project. The RMP describes how you will elicit, analyze, document and manage the project’s requirements process.
Program and Course Outcomes
This assignment is directly linked to the following key learning outcomes for the course:
• Know the general requirements collection process steps and become familiar with best practice requirements collection techniques
• Practice techniques for eliciting stakeholder’s input to develop effective and feasible requirements success criteria
• Distinguish between different types of requirements (functional, Non-functional, technical, project, regulatory, etc.)
• Analyze techniques used to differentiate the needs vs. wants of stakeholders for a project: MoSCoW technique
• Understand the purpose of an effective Requirements Management Plan (RMP): Create an RMP
Essential Components & Instructions
Requirements Management Plan (RMP) Overview
The Requirements Management Plan is a key document in managing the scope of your project. The RMP components describe how you will elicit, analyze, document and manage the requirements of the project. Specifically, the RMP will outline the up-front gathering of high-level project and product requirements, as well as the more detailed product requirements that you will collect during the project lifecycle.
Most importantly, adhering to an effective requirements management process helps the project team focus on the requirements that have been developed and maintains the integrity of the requirements throughout the lifecycle of the project.
Requirements Management Plan (RMP) Instructions
The following template outlines common RMP sections:
- Executive Summary: In 1-page, write this Executive Summary on the key components of your RMP.
- Project overview: In 2-3 paragraphs, briefly describe the purpose of the project for the readers that have not seen your Project Charter nor understands the project’s key business objectives.
- The requirements gathering process: In this section, describe the process that you will use to elicit, analyze and document the project requirements. Identify and describe at least 4 requirements collection tool(s) & techniques you intend to use to collect project requirements from key stakeholders. Moreover, please comment on collection method’s efficiently to collect necessary project requirements. (Note: project efficiency should focus on timeliness to collect requirements vs. cost to collect requirements)
- Key Roles and responsibilities: In this section, please list the roles & responsibilities of at least 5 key project stakeholders who will be involved with gathering, creating and managing the project requirements throughout the project ‘s lifecycle.
(For example: Roles could include the project manager, project sponsor, business analyst, customer(s) or project team SMEs, or other key stakeholders. Responsibilities could include requirements elicitation, change management, requirements creation, testing and approving requirements. etc.)
- Assess Requirements: From your project case overview and individual research, write at least 6 requirements. Please ensure requirements follow the SMART guideline for effective project requirements. Moreover, each requirement must be categorized based upon the project’s defined strategic objectives & goals.
Please categorized each requirement into Must Have, should have and Could have/Would like to have (pls refer to the MoSCoW Method). You can use a table for this section.
- Requirements Collection Timetable: Using at least 3 requirements collection methods, briefly identify a requirements collection schedule for your project. Please justify your timeline based upon key resources availability, their assumed skills and the collection method used.
- Requirements traceability: Please describe how the team will track and manage requirements from requirements elicitation to identifying project deliverable, WBS / project schedule development and scope change management. [Hint: Review 100% rule].
- RMP Approval: Identify who will review and approve the RMP. Please comment on how you will communicate the RMP to project team and key stakeholders.
Assignment Format
Below are some key guidelines you will want to ensure you follow in creating this assignment.
Think of this short list as a quality control checklist:
• Requirements Management Plan template must be complete, pls submit in MS Word (.doc, .docx) or .pdf formats
• You should format the document professionally
Project Scenario
Proposed: Enterprise Project Management (PM) Tool Selection and Implementation Project
You are part of a high-powered consulting group (EAA Inc.) who provides our clients with professional services on how to assess, recommend, develop and implement project management processes and systems. We have been in business for 10+ years – and we are considered the “Best and Brightest!”
I have just returned from a Board of Directors meeting with a large non-profit client (TestSmart Group [TSG]) who provides educational testing services under contract for Massachusetts State, Department of Education (DoE). I have been working with this client over the past few years to improve their organizational project management maturity.
For the first 3 years, our PM improvement focus and objectives with TSG has been on establishing formalized project management processes, simple tools & templates and significantly enhancing TSG’s PM capabilities for the following PM disciplines:
- Project Charter / Project Concept Documentation
- Requirements Elicitation Process & Project Scope Statement
- Work Breakdown Structure (WBS) and Project Schedule Development Processes
- Elaborative PM Estimating – Resource Effort and Task Duration
- Critical Path Scheduling and Management
- Project Communication Planning
- Integrated Project Scope Change Control Management Process
- Project Risk Management
- Project Issue Management
- Project Quality Management Planning
- Lessons Learned – for continuous project management maturity & process improvement
- Program Management
- Project Portfolio Life-Cycle Management (to include: Project intake, selection, prioritization and project sequencing)
Background: The work we have done with TestSmart Group (TSG) to date has included building the Project Management methodology frameworks and processes, conducting training for project managers and their project teams on effective project management, and coaching teams in the development of their project plans.
Moreover, we have most recently implemented processes and templates for program and project portfolio management to add strategic-level project governance to TSG. We have noticed that implementing enterprise wide project management tools and capabilities has a significant impact on TSG’s culture and how they provide services to their clients. Bottom-line is that implementing PM has been significantly changed the TSG organization.
Critically, we have realized that leading the organization change has been a big part of this long-term initiative. Moreover, we have also established an Enterprise PMO to be TSG’s ‘Center of PM Excellence!” Overall, we are excited about the improvement results and are actively preparing for next steps in our PM maturity journey with TSG – our most valued client.
Currently, TSG uses Microsoft SharePoint as a PM work product management repository and project team collaboration tool. Consequently, there is a project SharePoint site created for each project and teams have been using a number of different tools to manage projects including MS Project, Excel Gantt Charts, and Jira (www.atlassian.com/Jira).
While these tools are mostly working well, the organization has realized that they may be able to increase productivity and improve the quality of project planning and execution by using more current and productive PM collaborative tools. As a result, TSG has asked us to engage in a project to accomplish the following objectives [high-level project scope elements gathered by the initial business case done in 2016]:
- Define a process for assessing and selecting the “Best Value” automation tool that will meet all TSG’s PM requirements for needed PM processes, tools & templates and project planning and collaboration.
- Identify, procure, and then successfully implement a project management software tool(s) solution for TSG.
- Project work also must include the design & development of PM training for the selected PM automation tool.
- Providing PM automation tool familiarization training for TSG’s managers, all project sponsors, portfolio managers, program managers, project managers and project team members.
- Project must also identify and train advanced users (i.e. train-the-trainer) to serve as operational support for teams using the selected tool(s).
These were some high-level requirements and solution objectives gathered a few years ago. Obviously, TSG management expects us to fully gather additional requirements for this initiative. Therefore, requirements still need to be collected, prioritized and validated. TSG expects that the solution must support project managers at the project level, program managers in managing their enterprise programs and to provide the functionality for TSG’s senior managers to make informed decisions at the enterprise project portfolio level for project approval, prioritization and sequencing. Moreover, the selected solution must provide effective, accurate and timely project progress information and metric dashboards for TSG managers and TSG’s Board of Directors.
Currently, Microsoft Project (application) is used for scheduling single projects; however, it is not clear that the Microsoft Project Portfolio suite is the optimal solution for managing the integration of projects throughout the entire organization.
High-level proposed roadmap of action(s):
I had a high-level 1-hour meeting with TSH senior managers. At this meeting we identified some high- level actions that may help us move the project forward.
– TSG PM
stakeholder requirements need to be better defined, validated and categorized,
– Potential vendor solutions/tools need to be analyzed,
– Then an enterprise-wide PM automation solution selected and planning for implementation to be determined.
Notes I made to myself at this meeting:
- – The requirements / scope information here is really not fully adequate for an effective implementation project plan.
- – Therefore,
my consulting team must conduct [real-world] outside research to identify
current solution functionality & more detailed requirements to assess,
select and implement an effective enterprise wide PM solution.
[A good practice is to actually go to vendors’ website to see tool functionality and services).
- – Moreover, client’s current PM processes need to be considered and integrated into the chosen solution.
- – Don’t’ forget the extensive training that needs to be identified, designed and delivered – management is not talking about this!
- – Job aides & process guides, on how to use the chosen PM solution, must be developed for all TSG’s critical project management roles. I can’t think of them all??
- – Get with the team -so much to think about and begin Scoping!!
Action: At this meeting, I quickly volunteered our team to assess and create project requirements planning and scope development plan for this critical project(s). Taking points with consulting team:
- We have 3 weeks to begin work with stakeholders to define a more detailed level of clear and testable requirements.
- Once we better understand the project requirements, we have an additional 3 weeks to develop a viable Project Scope Statement (PSS) and detailed Work Breakdown Structure (WBS) to present to the Project Sponsor for review and approval (in Week 6).
- This is a very tight timeline so, to pull this project scope together, we need to quickly get to work!
I have provided preliminary information in the (very rough order of magnitude [ROM] – high level) draft scope, below.
I expect over the next few weeks the team must add additional insights to adequately complete this document for approval.
I am sure here are many things I have overlooked – team – please feel free to add additional sections to best define the scope of the project.
I (Professor) will play the role of Project Sponsor for this initiative.
Some High-Level Project Objectives/ Success Criteria (from old Project Charter):
- Define requirements to identify, procure and implement a project management software tool(s) for use within TSG.
- We have 6 weeks for the project scoping work – planning is to commence on 1st Qtr. – 2020 and project plan integration must be completed by summer 2020 for a Board decision.
- Since the work is budgeted for calendar years 2020 and 2021, the entire project needs to be complete by 12/31/2021.
- The organization’s leadership is committed to using this solution – beginning of 2022.
Some Project Charter elements already identified:
Assigned Project Manager: Each team select a PM
Sponsor (Charter): <Professor>, representing the non-profit client (TSG) will act as sponsor and approve all plans, budgets, requests for contingency and proposed project scope changes.
“Top Down” Project Scope Description: The work of this initial project is to define the scope for PM solution selection, procurement and plan to successfully implement an enterprise project management software tool(s) for use within TSG:
Initial Scope & Known Constraints – more to be determined by team:
|Work Includes (In Scope)||Work Does Not Include (Out of Scope)|
|At least 4 potential vendors must be evaluated. (Note: real world PM tool research is highly recommended)||Ongoing maintenance of the selected tool set.|
|Establish advanced users training / ongoing mentoring of TSG’s PMO on how to use the selected tool to improve TSG’s project management practices.||No infrastructure requirements need to be planned for – this is already being done by TSG’s IT Department.|
|Create and conduct training sessions for all TSG’s critical stakeholders.||Upgrades to hardware or connectivity. Also, security considerations have already been identified by TSG’s IT Dept.|
|Assess impact of new PM tool to TSG’s organization. Then, create a proposed effective organizational change management strategy initiative (follow on work).||Updates to existing project management templates, processes and methodology documentation.|
Some known Acceptance Criteria – more to be developed by team:
- Project, program, portfolio managers and team members are able to use selected tool in the execution of their project work. Success Metric: 100% of PMs, PgM and PPM team have tool access, and are fully trained on tools. Sponsor Approval / Sign-off required.
- Current TSG project management processes and templates are integrated with selected tool. Success Metric: 100% of all current PM tools, process guide, templates are successfully deployed into new PM tool. Sponsor Approval / sign-off required.
- New PM tool(s) are available for enterprise wide use by 1/01/2022. Success Metric(s): PM tool functionality approved by all stakeholders, tools integration and systems testing completed. Sponsor Approval / sign-off required.
- Organizational Change Management plan proposed. Success Metric: Proposed TSG Change Mgt Plan reviewed and approved by project sponsor for presentation to Board.
Project Exclusions – see out of scope above.
Project Constraints
- PM solution selection Budget – $650,000 (USD) has been approved for project.
- Project Sponsor holds an additional $100,000 in management reserves – Sponsor controls this contingency budget. TSG has an additional $500,000 in reserves. Only TSG’s Board of Directors can approve and allocate these funds.
- Quality – All customer requirements identified as ‘deliverables’ and adequately planned for enterprise wide implementation and deployment. Project sponsor reviews and approvals are signed off before the project is declared complete. Project teams may leverage PMI’s PMBOK processes as reference standards.
Known Project Assumptions
- No authorized project work beyond project scope & planning
- Core project team (you) will be hired as consultants for execution of the project
- Client’s (TSG’s) human resources are also available to be SME’s and to be assigned to the project.
- The core project team (you – consultants) may only devote 75 percent of their schedule to this project (Standard: 40 hours/week). Overtime is available and must be approved by project sponsor.
Minimal Major Milestones (of course I would expect to see more)
- Project Start Milestone • Milestones for each Project Phase (Project Initiation, Project Planning, Project Execution and Project Shutdown)
- Tool Selection Completed Milestone
- Project Planning Complete and Approved
- Entire Project Completed
- Other Milestones – To be determined by Project Team
Follow on Project Scenario Details
Periodically, I may provide additional information to teams – like recommended project scope change. This additional information may not be the same for all teams. The goal here is to have each team react to ‘real world’ project issues and assess impact to the project.
Do You Know That our Professional Writers are on Stand-by to Provide you with the Most Authentic Custom Paper. Order with us Today and Enjoy an Irresistible Discount! | https://elitepapershub.com/2020/05/27/roles-responsibilities-of-key-project-stakeholders/ |
For sale recently, restored farmhouse in Cortona in Tuscany. On the border between Umbria and Tuscany, near the beautiful Cortona, for sale farmhouse completely restored in 2008, with 5 bedrooms, 5 bathrooms and in addiction 2500 sqm park with swimming pool. The Location is panoramic, with Trasimeno Lakes view and surrounding land, and with privacy but not isolated, in a few minutes you can reach many little town with all the services. The building of 250 square meters, on two levels, is composed as follows: Ground Floor - Entrance, living room with fireplace, large kitchen with dining area from which you can access the garden through a French door, bedroom, study, bathroom and closet. On the first floor- reachable from both internal and external staircase with loggia, there is a large hallway, 4 bedrooms and 4 bathrooms. The structure has been restored using only the highest quality materials, the rooms are large, well-lit and equipped with air conditioning, there is also the possibility of creating an outdoor pergola. Energy Class "F", Ipe 165.915 kWh / m2. The land, included in the sale, has an extension of 2500 square meters entirely used as a park with a pool of 5 * 10 meters; everything is fenced and equipped with an automatic irrigation system. Distances: In 20 minutes (15 km) you can reach the centre of Cortona, Montepulciano, Castiglione del Lago and Foiano della Chiana. The nearest airports are Perugia (46 km, 40 min.) and Florence (110 km, 1 h 15 min.). | https://www.realestate.com.au/international/it/cortona-tuscany-120030326495/ |
A database of pairwise, structure-based alignments for structurally analogous motifs in proteins. Applications of this database may range from protein evolution studies, e.g. development of remote homology inference tools and discriminators between homologs and analogs, to protein-folding research, since in the absence of evolutionary reasons, similarity between proteins is caused by structural and folding constraints.
1 - 7 of 7 results
tune Filters / Sort by
1 - 7 of 7 results
Compiles structural neighbors of proteins deposited in the Protein Data Bank (PDB). FSN allows to search by a PDB code as a query and provides a list of structurally similar proteins sorted by increasing P-values, as well as links to various statistics about this group of proteins. The database provides a table lists structurally similar proteins with each individual protein in a single row, providing detailed information about each comparison and links to PDB and detailed FATCAT result page.
Provides analysis of protein topology and its modular architecture. ProLego offers an alternative approach to study protein structure topology. It compiles an extensive topology database analyzing different sets of non-redundant representative protein datasets. It can be used for identifying constituent topological modules in proteins of interest, which could be used as “lego-blocks” in protein designing.
Provides access to data from studies on direct coupling analysis (DCA) ability to differentiate between properly folded and misfolded structures. DCA vs. Misfolds is dedicated to studies on DCA, a method that can assist researchers in studying protein structure.
Contains incorrect conformations to improve protein structure prediction. Decoy ‘R’ Us provides a resource that allows scoring functions to be improved. It can be used to evaluate and improve scoring function performance for predicting structure, to elucidate the physical nature of protein–protein interactions or to assess the degree to which biologically relevant functional sites are preserved in predicted structures.
Combines structural coverage leveraged from homology models and experimental protein structures. PMP uses a 3D molecular models in biomedical research to find both experimental structures and theoretical models for a given protein. PMP is an open project for the community offering a unique interface to visualize structural coverage and to analyze the variability of a protein.
Compares predictions of several fold-recognition techniques to the Saccharomyces cerevisiae genome. SPrCY is a database which allows users to search, browse and analyze the generated predictions. It offers an interest with the computational biology community, and the new structural and functional annotations for the yeast genome to help guide new experimental research on this important model organism. | https://omictools.com/protein-structure-comparison-data-category |
Wild African elephants are voracious eaters, consuming 180 g of food per minute. One of their methods for eating at this speed is to sweep food into a pile and then pick it up. In this combined experimental and theoretical study, we elucidate the elephant’s unique method of picking up a pile of food by compressing it with its trunk. To grab the smallest food items, the elephant forms a joint in its trunk, creating a pillar up to 11 cm tall that it uses to push down on food. Using a force sensor, we show the elephant applies greater force to smaller food pieces, in a manner that is required to solidify the particles into a lump solid, as calculated by Weibullian statistics. Elephants increase the height of the pillar with the force required, achieving up to 28% of the applied force using the self-weight of the pillar alone. This work shows that elephants are capable of modulating the force they apply to granular materials, taking advantage of their transition from fluid to solid. In the future, heavy robotic manipulators may also form joints to compress and lift objects together.
1. Introduction
Wild elephants browse and graze for up to 18 h per day [1,2], consuming over 200 kg of vegetation per day [3]. Thus, on average, an elephant eats 180 g of food, or the weight of two corn cobs, per minute. Even in captivity, elephants (figure 1) continue to consume food at up to half this rate [3]. To eat at these high rates, an elephant uses its trunk to pick up as much food as possible each time it reaches out. This behaviour is analogous to using a fork to pick up as many noodles as possible before each bite. Picking up multiple objects at once requires practice and physical intuition as to how piles of materials behave under applied forces [4,5]. While little is known how elephants perform this feat, there is a growing interest in robotics in conducting similar tasks. For robotic manipulators to work in the real world, they will have to deal with multiple objects in cluttered and unpredictable environments [6,7]. The goal of this study is to elucidate how elephants manipulate multiple objects at once. Figure 1. The indoor enclosure where experiments are conducted. During experiments, the elephant turns to face the force plate and video cameras and protrudes its trunk through the enclosure. (Online version in colour.)
We focus here on piles of granular materials, collections of discrete, solid, macroscopic particles. Examples include construction material, such as sand and gravel, as well as food items, such as flour and chia seeds. Sand and gravel are often pushed with bulldozers and grabbed using construction cranes with an end attachment called a clamshell [8]. In both cases, a dustpan-like device is slid underneath the pile of materials in order to lift it up. Elephants use a different mechanism: they squeeze the particles together, jamming the grains which cause the pile to solidify. Such a mechanism might be used to help soft robotic grippers to pick up multiple objects together [9–12].
The elephant trunk is similar to other boneless organs in nature such as the octopus arm, and the human tongue and heart. These organs are composed of a tightly packed array of muscle and connective tissues. They are known as muscular hydrostats and are composed of interdigitated muscle fibres arranged in three dimensions. They thus lack the discrete muscles of rigid skeletal support systems [13]. The elephant trunk is the largest muscular hydrostat on land, making it subject to substantial gravitational forces.
In this study, we investigate the behaviours used by elephants to pick up multiple items simultaneously. We begin in §2 with our experimental methods for filming and measuring the forces applied by elephants. We proceed in §3 with our mathematical models for the squeezing force applied to the food and the granular physics of jamming. In §4, we present our experimental results, focusing on the forces applied to pick up different sized food items. In §5, we discuss the implications of our work and suggest directions for future research, and in §6 we state our conclusions.
2. Material and methods
2.1. Elephant training and husbandry
All experiments are performed on a 34-year-old female African elephant Loxodonta africana over several weeks in the summer of 2017. The elephants are trained to perform a number of routines for visitors, including a demonstration of basic movements of the body, and reaching for food using the trunk. All experiments are supervised by the staff at Zoo Atlanta.
2.2. Measuring trunk density and trunk weight
Using an elephant trunk that is cut into four sections, all of which are stored in a freezer with a temperature of − 20°C, we are able to collect length and mass data. While the trunk is in the shape of a frustrum, the last 23 cm can be approximated as a hollow cylinder using the equation:
2.1
trunk
0
v
0
trunk
steak
whereis the density of the trunk, andis the trunk outer radius,the height, andthe inner radius of each of two nostrils. We measure the mass of the frozen trunk section as= 2.35 kg and its height= 23 cm. The trunk section has an outer radiusof 52 mm, and an inner radiusof 15 mm. Thus the volume of the trunk can be calculated as. Using equation (2.1), and the weight of the trunk section, we calculate the average density of the trunk tip as= 1.5 g cm. This value is above the density of lean boneless cow muscle,= 1.2 g cm, possibly because of desiccation in the sample [ 14 ].
To estimate the weight of the trunk, we photograph the elephant when its trunk is in a relaxed position (electronic supplementary material, figure S1). We measure by hand the tip diameter d 1 = 12 cm, and so infer from the photograph that the trunk has a length L trunk = 1.9 m and is widest proximally, with a diameter of d 2 = 38 cm. Approximating the trunk as a frustrum with two nostrils, its volume is . The total mass is m trunk = ρ trunk V frustrum ≈ 150 kg (see details in the electronic supplementary material).
2.3. Grabbing force measurement
To prepare food for the elephant, we cut by hand rutabaga and carrot into cubes of side length 10 mm, 16 mm and 32 mm. We also scoop wheat bran with grains of characteristic size L ∼ 2.0 ± 0.5 mm, and volume V ∼ L3 = 0.008 cm3. The food is arranged by hand into a small pile in the centre of a force plate (Accugait, AMTI, USA) for each trial. We separate the food into piles of approximately the same size: this means 50 g of bran and 100 g of cubes in sizes of 10 mm, 16 mm and 32 mm. Since wheat bran has a density of ρ = 0.17−0.25 g cm−3, then M = 50 g of bran has approximately N = M/(ρV) ∼ 40 000 particles in it. Thus, the number of particles that we test varies over four orders of magnitude, from four particles to 40 000.
Figure 1 shows the location where experiments are conducted. The elephant stands behind the bars of an indoor enclosure and extends its trunk through the bars to reach food. Food is placed on the force plate, whose edge is a horizontal distance of 46 cm from the enclosure. Two video cameras (Sony Handycam, Japan) are placed in the bird’s-eye view and side view of the force plate. An indicator light (Massimo Retro LED, USA) activated by remote control is used to synchronize the force plate and cameras.
We start every experiment in the morning at 9.30 EST and finish it within an hour. First, the force plate, indicator light and cameras are installed. The force plate is zeroed and the indicator light is turned on. Each trial begins by the curator instructing the elephant to retrieve the food. The elephant draws close to the force plate and stretches its trunk to grab, as shown in figure 2. The two cameras start to record the scene and the indicator light is turned off to synchronize both cameras. The real-time contact force data are captured at the same time. For each of these food sizes, we conduct six trials, providing a total of 24 trials, of which 16 were analysed. The remaining eight trials were not analysed because the elephant performed a trunk wrapping rather than a jamming motion. The rest time between experiments is about 2 min. Figure 2. Time sequence of elephant trunk sweeping and grabbing a pile of wheat bran. (a) The trunk locates the force plate. (b) The trunk tip sweeps for about 5 s to compact the bran. (c) The trunk tip pushes downward to jam the bran using both finger-like extensions on the trunk tip. (d) The trunk detaches from the force plate, carrying food to the mouth. (Online version in colour.)
2.4. Image analysis to locate the trunk joint
The location of the joint is found using image analysis, which is discussed in detail in the electronic supplementary material. We begin with a guess, by estimating by eye the location of the joint, defined as the point at which the elephant begins to form the distal end of its trunk into a distal pillar. We binarize the image using Matlab and then use image analysis tools to extract the points characterizing the most distal and proximal ends shown in the image (electronic supplementary material, figures S2–S6). Two lines, shown in red dashed lines in figure 3, are fit to each of these series of points, and their point of intersection is calculated. This intersection point is the new location of the joint, and it often falls quite close to the initial guess. The coordinates of the joint are used to measure the height of the trunk pillar. In figure 3, the joint is shown by the white point and the height of the pillar by the yellow dotted line. In the next section, we present our mathematical modelling tools which we use to rationalize the shape that the elephant trunk takes to grab each object. Figure 3. Trunk configuration when jamming food for (a) 32 mm cubes, (b) 16 mm cubes and (c) 10 mm cubes and (d) bran granules of diameter 2 mm. Note the carrot cubes are orange and the rutabaga cubes are red. The red dashed line is tangential to the top 50% of the trunk above the joint. Note that the trunk is straight when grabbing cubes with a side length of 32 mm, but then forms a joint when grabbing smaller pieces. When grabbing bran, the vertical part is the longest, reaching up to 11 cm. (Online version in colour.)
3. Mathematical modelling
3.1. Forces by the trunk pillar
To pick up granular materials, horizontal squeezing forces must be applied to the pile. While humans can use two hands to squeeze the pile, elephants are constrained by the anatomy of their appendage. African elephants like the ones in our study have two finger-like appendages at the tip of their trunk. These fingers push on food as shown in figure 4a,b. Because the fingers are oriented at a non-zero angle α relative to the vertical, it allows the elephant to transduce downward forces into horizontal forces that contract the pile together allowing it to be picked up. The idea is similar to scooping flour up from a table by squeezing it between the fingers and palm of one hand. Figure 4. Schematic of forces applied to the force plate and to the food. (a) Schematic of the elephant trunk, with pillar weight m v g, applied force at the joint of F m . The force plate responds with a force F plate . (b) Schematic of forces applied to the food pile. (c) The food is able to lift off the ground because of lateral forces F x . Here is a jammed arch of granular particles due to the application of the horizontal force of F x . (Online version in colour.)
Rather than consider the mechanics of the food–finger interaction, we consider a force balance on the entire trunk pillar as control volume as shown figure 4a. Forces arise from the following three components: applied force F m , the plate’s reaction force F plate and the weight of the trunk itself m v g where m v is the mass of the trunk and g is gravity. The vertical force balance may be written as
3.1
plate
m
v
3.2. Mathematical model of jamming force
In other words,: the force on the force plate is equal to the self-weight of the pillar plus any forces the elephants apply. We proceed by presenting a model for the force required to solidify the food particles.
We create an ansatz model interpretation of these results that take into account the fundamental granular nature of the food. Unlike continuum solids, force is propagated through granular materials in discrete chains, which can be deflected due to oblique particle contacts. In order for a collection of granular food to be lifted as a solid, a stable arch must span the entire two-dimensional area at the base and be of sufficient strength to withstand the weight of the particles above. The statistics of arch formation have been studied in two-dimensional [15] and three-dimensional [16] hoppers, with two important findings. First, the relevant length scale is the particle size; this means that spanning arches in small foods can be thought of as ‘longer’ than those in larger foods in that they span more particles. Second, weakest link theory explains the intuitive finding that longer arches are weaker, and therefore less common, than shorter arches. The various statistics involved in arch formation and destruction can be quantified through random mean-field approximations of particle location. Here we apply a related Weibullian weakest link analysis [17], to rationalize why the elephant applies greater forces to pick up smaller food particles.
Weakest link statistics were developed in 1939 by Weibull to explain the strength of continuum materials [18]. The analysis builds on the single assumption that a long sample comprises many smaller elements that are statistically independent. It was subsequently [17] applied to explain the stick–slip yielding of geometrically cohesive granular materials under extensional strain. Here, the key idea is that each particle contact has a probability of failing that is independent of the state of the other particles. In order for an arch or pile to be stable, all contacts must independently be stable. A review of the entire model, including its extensions to identify failure location and time-dependent failure can be found in the previous literature [19].
We relate by analogy the failure probability with the force needed to prevent an arch break-up. Granular contacts must support vertical forces at least equal to the pile weight. Consider four grains forming an arch supporting the weight of a particle above, as in figure 4c. As the angle becomes more oblique, the normal force required to sustain the weight diverges as 1/sinθ. In real grains, the normal force is supplemented with a frictional force, both of which increase with confining (lateral) forces applied by the elephant. In this interpretation, if the force is not large enough to maintain the arch, then the elephant can re-establish stability by applying a larger force, i.e. weaker piles are stabilized by larger applied forces. We now show that the mean force at failure decreases as a power law with the particle size, S.
Following Weibull’s original weakest-link analysis, we first assume that for small values of applied force F, the probability of a differential length δL to fail depends linearly on δL and increases with applied force F as some undetermined power law, δY ≡ FmδL; the probability for that differential length to not fail is 1 − δY. For a longer sample composed of i multiple units to not fail, each individual sub-unit must not fail. The probability that all units simultaneously will not fail is
3.2
3.3
3.4
3.5
3.6
3.7
Figure 5. The relationship between applied force and food size S. The solid points represent the applied force recorded by the force plate, and the open points represent the self-weight of the trunk pillar, defined as the region below the joint. The trend line is the power-law fitting of the applied forces. The inset shows the relationship between the height H of the trunk pillar and food size S. Error bars show the standard deviation of the measurement. The self-weight of the pillar has a standard deviation of 0.5 N, which is too small to see in the graph; for that reason, the inset is shown. (Online version in colour.)
4. Results
whereis a constant introduced for dimensional reasons and the product is over allunits. We assume that the probability of an individual unit yielding is small compared to 1, in which case we can make the approximationThe sum of the differential lengths is just the total sample length,. In experiments, the elephant trunk is picking up samples of approximately the same length; as the food sizedecreases, more particles are needed to span that space. The indexincreases as 1/, and so we can rewrite equation (3.2), which gives the probability that the chain will not break, in terms of the applied forceand the grain sizeasThe probability that the chain of particles of sizewill fail at the applied forceis then equation (3.4) subtracted from unity, orWe now reproduce the calculations from [ 17 ] to find the average yield force as a function of grain size. Equation (3.5) is the probability that a chain of particles of sizewill fail at the applied force. Piles are not subjected to instantaneous forces, however. For failure at aforceto be observed, the pile must not fail at the lesser forces applied. The total probability to observe failure at forceand sizeis therefore the product of equations (3.4) and (3.5):where the prefactoris included for normalization so that. The mean force observedis thenis the average force required to break a chain of particles of size. Our hypothesis is that when the pile weight exceeds this force, the elephant stabilizes the pile by applying a larger lateral force that increases the friction forces within the pile. For simplicity, we assume that the applied forceneeded is first order inversely proportional to the pile strength, and so we arrive at the conclusion thatshown in figure 5
We filmed 24 trials of the elephant grabbing food. In a third of the trials, the elephant curled its trunk around the food to grip and lift it. Figure 6 and electronic supplementary material, video S2 shows the elephant curling to grab an entire pile of 10 mm cubes. This technique is successful at obtaining more than 80% of the food items. The remaining 20% of food items are fetched on a return trip of the trunk. Each curling action takes 6 ± 2 s (N = 4). For the remainder of this paper, we focus on the elephant’s most typical method for grabbing piles of particles: formation of a joint and downward pushing to jam the particles. Figure 6. Food grabbing by curling the trunk. (a) The trunk curls and sweeps the food together. (b) The trunk carefully squeezes the food in a loop to carry it to the mouth. (c) The trunk loop holds the food. (d) The food is picked up by the trunk. (Online version in colour.)
The elephant’s method of grabbing bran is shown in figure 2 and electronic supplementary material, video S1. The elephant first extends her trunk to locate the force platform. When this occurs, the food pile is usually missed by 10 cm, which is suggestive of the elephant’s poor vision. Once contact is made with the platform, the elephant sweeps the food into a pile with the tip of her trunk. During the sweeping process, she appears to keep the trunk oriented diagonally, aimed directly toward the food. However, grabbing the food requires substantial horizontal forces to stabilize the particles. Thus, once sweeping ceases, she pushes downward while spreading her trunk’s two finger-like extensions, as shown in figure 4. Then the food is taken into the mouth by curling the trunk (figure 1).
This sequence of events corresponds to changes in the applied force, which we measure with a force platform synchronized to the video. The time course of the contact force is shown in figure 7, for 32 mm cubes, 16 mm cubes and bran. When the elephant trunk makes the first contact with the plate, a force peak of magnitude 20−40 N is reached for a fraction of a second, associated with impact of the trunk with the plate. We believe this force is large because the trunk is heavy; by tracking the trunk tip when it approaches the force plate, we find that the elephant trunk actually slows in speed by before impact with the plate (electronic supplementary material, figures S7–S10). This reduction in speed suggests that the elephant can anticipate the position of the force plate. After the initial impact, the elephant rests part of her trunk on the scale as she sweeps the food, showing a plateau in force of 10–20 N for a duration of 4–10 s. This force is likely required to ensure adequate contact with the scale to perform the sweeping action. The contact force doubles to 30–40 N for a second when the elephant picks up all objects, except for the 32 mm cubes. In figure 7a, the force applied when picking up the 32 mm cubes is only 7 N. By watching our videos synchronized to the force platform, we observed that the onset of peak contact force coincided with the elephant bending its trunk from a straight configuration to one with a kink, or joint. Figure 7. Real-time contact force while grabbing (a) 32 mm cubes, (b) 16 mm cubes, and (c) wheat bran (2 mm in diameter). A force peak of 20–40 N is made when the elephant trunk first reaches the force plate. The force then drops and forms a plateau of 10–20 N. At the moment when the elephant pushes down the food to pick it up, the contact force doubles to 20–40 N. (Online version in colour.)
Figure 3 shows the configuration of the trunk for each food item, at the point where the applied force is highest. The elephant forms joints in all 12 trials except for trials involving the largest food size, 32 mm cubes (figure 3a). We characterized the pillar by a height H, shown by the dotted yellow line in figure 3. When picking up bran, the trunk pillar has a height of 11 ± 0.39 cm. In comparison, when picking up 16 mm cubes, the elephant uses a pillar height that is one third as tall, of height 3.9 ± 0.55 cm. Clearly, the elephant has a great deal of control of the height of this pillar. Using the density of a deceased elephant’s trunk, we calculate using equation (2.1) in §3.1 the weight of the pillar.
Figure 5 shows the pillar weight (open points) and applied force by the elephant (closed points) as a function of the different food size S. As objects decrease in size, the pillar weight and applied force both increase. When picking up the smallest object, the elephant applies a force of 48 ± 2.1 N and generates a pillar of 11 ± 0.38 N in weight. For all the food items except for the 32 mm cubes, the pillar weight is 20–30% of the force applied. The elephant does not form a pillar, for the 32 mm cubes, but still applies a force of . We speculate it is the minimum force resolution that the elephants can sense.
The black line in figure 5 is a power law fit to the applied force by the elephant, whose equation is given by
4.1
5. Discussion
whereis the food size. Equation (4.1) is a good fit to the experimental measurements (= 0.76). While we cannot predict the exponent nor the prefactor in equation (4.1), the theory in the math modelling section correctly predicts that the exponent has a negative sign. The bran of size 2 mm requires nearly 50 N of force, more than three times the force of the 32 mm carrot cubes. Why do smaller objects require more force to pick up? The difference between the pile of carrot cubes and bran is that the bran involves a far greater number of particles. When the particles are squeezed together to be picked up, each particle has a small chance of failure. Thus, the bran pile requires more applied force to overcome the accumulated failure probability of the large number of grains involved.
Although the elephant trunk lacks bones, the formation of a joint mimics a common vertebrate strategy to reach out and grab objects. The human upper limb, for example, has seven degrees of freedom. These degrees of freedom make it possible to reach out into arbitrary points in three-dimensional space and grab objects, as well as perform twisting motions in all three directions. An animal with more joints has more degrees of freedom to accomplish tasks. But these joints also provide challenges too, as the animal must search through more potential solutions. This is why appendages without bones, such as the elephant trunk and octopus arm, have both demonstrated the formation of joints. The octopus forms a joint like the elbow only when retrieving food [20,21]. Our study shows that the use of joints might be more common than once thought.
In our study of captive elephants, we prepared cubic food items that the elephant would never find in nature. Nevertheless, wild elephants may still apply the strategies we observed if they need to press downward with their trunk while feeding. Wild elephants eat grasses, small plants, bushes, fruit, twigs, tree bark and roots. To remove the bark from a tree, vertical forces are required, and its possible the elephant may form joints for this task. Now that we have observed the formation of joints, future work will determine how often elephants use this strategy.
Long flexible robots have long been of interest to the robotics community. Such researchers have turned to snakes, octopus and elephants for inspiration. However, even among these animals, the elephant stands out because the trunk can apply the greatest forces. For elephant-inspired robots to apply large forces, they will inevitably become larger. We surveyed four elephant-inspired robots whose weights were reported [6,22,23]. On average, their weight is 5 kg, which is nowhere near the elephant trunk. Nevertheless, a number of elephant-inspired robots have sufficient degrees of freedom that they could be used to generate joints [6,11,24]. In particular, the elephant-inspired robot by Mcmahan et al. [7] can perform many of the corresponding motions observed in our work. For example, when this robot lifts an aluminium can, it cradles the can by forming a kink in its trunk, clearly showing that elephant robots have the ability to form joints (see video accompanying the paper) [12].
In our study, we observed the elephant is applying up to 47 N of force in order to pick up the 50-g pile of wheat bran. This means that the elephant must exert 100 times the weight of the pile in order to pick it up. We identified in this study that the weight of the trunk pillar provides up to 28% of the applied force. The remainder of the forces may also come from self-weight of the remainder of the trunk. The entire trunk weighs about 150 kg, or 1472 N. Thus, by simply relaxing just 3% of the weight of the trunk, it might generate enough force to compress the wheat bran. The entire trunk is about 1.9 m long when it is relaxed (electronic supplementary material, figure S1) [25]. We estimate that the 47 N of applied force would require the distal 46 cm of the trunk to be recruited to apply self-weight. In our experiments, the elephant was a horizontal distance of 46 cm away from the force plate, which provided a large constraint to the elephant’s grabbing. If the elephant were closer, it might generate a taller trunk pillar to help itself.
6. Conclusion
In this study, we investigate how elephants pick up piles of objects. The challenge in performing this task is that compressive forces must be applied to the objects so that they do not slip away. Using mathematical models, we showed that the greater the number of objects, the more compressive force must be applied. We test this idea in our experiments by providing elephants with food items varying from four to 40 000 in number. Elephants accordingly can vary the forces they apply by a factor of four, from 7 to 47 N. Using synchronized force platforms and video cameras, we show that the application of this force is accompanied by the formation of a kink or joint in the elephant trunk. The distal end of the trunk forms a pillar which provides up to 28% of the applied force. Forming joints may help reduce the energy required to reach for and grab food items, a task they perform for 18 h every day. The joint formation may also have application in elephant-inspired robots.
Ethics
All experiments are approved by Zoo Atlanta’s Scientific Research Committee and the Georgia Tech Institutional Animal Care and Use Committee.
Data accessibility
This article has no additional data.
Competing interests
We declare we have no competing interests.
Funding
This work was supported by the US Army Research Laboratory and the US Army Research Office Mechanical Sciences Division, Complex Dynamics and Systems Program, under contract no. W911NF-12-R-0011.
Acknowledgements We thank N. Elgart and Zoo Atlanta curators for helping us conduct the elephant’s grabbing experiments. We thank J. Reidenberg from Mount Sinai School of Medicine, New York for providing the frozen trunk section. We thank the Division of Mammals, National Museum of Natural History, and the Smithsonian Institute for their assistance with the frozen trunk section. We thank Dr Dalen Agnew for working with the frozen trunk section. We thank Sara Ha, Gerina Kim and Dhanusha Subramani for their early contributions.
Footnotes
Electronic supplementary material is available online at https://dx.doi.org/10.6084/m9.figshare.c.4257323.
| |
Hello everyone, it’s Cosmo the Library Cat. When I first started visiting the Grand County Public Library, I observed people sitting around a table and shuffling little colored pieces of cardboard around. This was strange behavior (even for humans!), but then I found out that they were working on jigsaw puzzles!
Apparently, people like to fit puzzle pieces together and then fit even more pieces together until they’ve created a complete picture. Personally, I think it might be more fun to chew on the puzzle pieces or bat them around but it seemed like the people working on the puzzles really enjoyed them. We don’t have a puzzle table out at the library right now, but if you have a library card, you can check out jigsaw puzzles and take them home to work on whenever you’d like. There are dozens to choose from, ranging from 500 to 1,000 pieces. My favorite puzzle at the library is called “Frederick the Literate” based on a painting by Charles Wysoki which shows a tabby cat blissfully napping amongst books on a shelf.
To borrow a library puzzle, you can browse online at catalog.moablibrary.org under the category “puzzles,” place a hold and use our curbside pickup service. Or if you come in for a personal visit, you’ll find the puzzles shelved with art books in our adult nonfiction section. Plus, the librarians are always happy to help you find whatever you’re looking for. So go ahead and take a puzzle home and put it together. Just be sure not to chew on the pieces, however tempting they may be. That way, someone else can put the same puzzle together later! Meow for now! | https://moabsunnews.com/2022/02/11/cosmos-corner-its-puzzling/ |
(1) The spirit and the matter being inseparable, the science of Scriptures and the science of the universe are together the one and only science. The prophets knew it. That is why, and although they had written the Book over thousands of years ago, they already spoke about the globe of the Earth, or about the circle of the Earth to show that it is round; of the atom of dust, indicating that there is something smaller there than what we see; the stars and their celestial bodies, meaning that stars have celestial bodies as the Sun has its own, and that our world isn’t an exception. They also speak about the Wheel filled with eyes inside and outside, showing the Galaxy completely inhabited; the chains of the pleiades Orion's belt, revealing the material ties of celestial bodies; and other things which still testify of their deep knowledge.
Space and the galaxies
(2) It isn’t indispensable to observe what the universe contains with telescopes, because it can’t be anticipated and demonstrated only by its material side. To grasp it completely, it is necessary to understand at first that it is a volume completely composed of matter; that it is the opposite of emptiness. Indeed, just as time, space exists only by the matter which forms it. That is why I said that between celestial bodies and between galaxies, space is a volume entirely of ethereal matter. We are certain that it is so, because a volume isn’t understood by forms or by limits, but only by the matter which composes it and within which there are distances, intervals. It is therefore very easy to see that the intersidereal and intergalactic space is a completely material volume, if only for the distances which separate celestial bodies and galaxies. But many other things will show us that space is completely formed by ethereal and subtle matter at the origin of the particles which composes celestial bodies, and these last ones the galaxies.
(3) So that there are no confusions in your minds, know forever that the galaxies which fill space here and there are wheels composed of stars with their celestial bodies, and clouds of vapors and dusts resulting from their work. With the solar family, we are inside one of these wheels and amongst the worlds which it contains and which are also as numerous as the stars. These worlds are the reason for being of galaxies which fill space.
(4) To know what is taking place within the galaxies, let’s give here some brief explanations on their composition. Let’s know firstly that, as beings are renewed, the celestial bodies which compose a galaxy are renewed likewise. That is why a galaxy is the biggest body composed of matter which integrates and disintegrates permanently, by the electromagnetic activity of the celestial bodies which are magnetized masses. We shall see that a celestial body is a sphere of which the metal part occupies almost all the volume, and that this sphere is magnetized by the celestial body from which it comes.
(5) Magnetization is an activity, a movement of matter which forms initially the magnetosphere of the celestial body, then the electrons which constitute the lines of force and the rings surrounding this celestial body. The magnetosphere is the essence of space which descends on the celestial body by condensing and by applying pressure on all the bodies. This intake of essence causes then the forming of the lines of force which leave from an hemisphere and arrive on the other one after forming rings in space. These invisible rings surround the celestial body here and there, perpendicular to the equator. And it is them that eventually give birth to other magnetized spheres, smaller, which are their satellites. These grow in their turn, and become full-fledged celestial bodies. A child will be able to grasp these explanations.
(6) At the moment, remember that the celestial bodies which compose a galaxy are magnetized masses having each a magnetosphere. These magnetospheres, which are felt very far in space, are added to one another to form a single and immense one which surrounds the galaxy. Conversely to that, the stars of that galaxy burn and are consumed by returning to space the essence which gave body to them. It is what constitutes the solar wind and the wind of all the stars which form together a very big galactic breath and an immense light. That said, the galaxies are surrounded at the same time by a magnetosphere (which is comparable to an inhalation) and by a breath (which is comparable to an expiration), because the magnetosphere is the essence which arrives on it, and the breath is the essence which leaves it. The magnetosphere is the intake of matter, which is a part of the INTEGRATION. The breath is the consumption of matter, which is a part of the DISINTEGRATION. Continual intake and consumption of matter give existence to the galaxies which, in this way, are perpetually renewed.
Birth and forms of the galaxies
(7) As the particle is born from another particle and the celestial body of another celestial body, the galaxy is born likewise from another galaxy. It is initially about a cluster of stars which, such as an embryo, forms within the wheel, then leaves it and develops afterwards. When a cluster gets loose from the branch of a galaxy which gives it birth, it has the shape of a ball which turns on its axis, and which takes a flattened shape little by little. This rotation movement, which the small galaxy acquires, is due to the stars which fade in its center; because the big void caused by the sudden disappearance of the immense magnetosphere of a star which fade, is immediately filled by the surrounding stars. It forms an inhalation, a driving force which pulls the stars towards this region where they disappear when their mass is completely consumed. A star consumes itself, and the moment comes when it is completely consumed. But the duration it took to disintegrate is infinitely longer than the one needed to be formed.
(8) All which is said will be put to observation. For the moment, let’s note that it is the disappearance of the huge and old stars that pulls the other stars towards the heart of the Wheel. This movement also makes the Galaxy turn on itself and forms its branches which coil in spiral. And by the centrifugal force which it exerts, this rotation also creates the separation of stars of different inertia. Because of that, we have to represent to ourselves the families of heavy masses (such as the solar family) in the lower third, the moderately heavy families in the middle third, and the light families in the center’s third, forming the bulb of the Galaxy. It is so, because a stellar family gets lighter in mass little by little as its planets destined to shine leaves it. The solar family will know this progressive lightening, when Neptune, then Uranus, then Saturn, then Jupiter will become stars and will leave it each in turn.
(9) Here is our Galaxy, seen from the side. This image makes us aware that such magnificence can’t exist in space without reason, aimlessly, without will and without purpose. However, as big as it is, it's only a luminous little point amongst the myriads of myriads of similar points distributed in the unlimited space. But it would appear such as we see it here if we went out of it and if we looked at it from the outside. Let’s imagine that we went away from it with the Earth, until it appears to us on the horizon. It is so vast that it is necessary to turn our head to the left then to the right to see its extremities. Thus its magnificence appears to us.
(10) From the edge up to the Sun, it’s only a small distance in which we find the worlds of the animal kind, as those who preceded us on Earth, and that of men. And from the Sun up to the center, there is, in this vast distance, the angels’ worlds in which we enter with the solar family which goes to the center of the Wheel. But everywhere our sights go, there are lands with greenery, seas and beings. Because in this immense house, there is not a single place where there are no living beings, no stars which shine in vain and pointlessly.
The depths of the universe
(11) If we could see our Galaxy laid over the horizon, as if we had gone out of it with the Earth and as if we looked at it from afar, wouldn't that be the biggest and the most beautiful spectacle which we would attend? Imagine then that we went alone even further in intergalactic space with the Earth, until we were halfway from the neighboring galaxy towards which we are going. From this place, this galaxy and ours are the same size and very small for us. We notice that we don’t distinguish their stars anymore. We see only the general light of both galaxies and not their stars which we can’t distinguish anymore. Everywhere we look in space, we can make out similar luminous points. These points aren’t stars but galaxies, myriads of galaxies in all directions and here and there which occupy all the volume of the universe. What a spectacle!
(12) Now let’s get away unlimitedly from our two galaxies, and let’s go to meet the following ones as fast as we would make it by walking on the stones of a ford. Let’s move forward! Let’s move forward! And let’s do it for a thousand of our years! During all this time, we never meet the end, and never are we in the nothingness where there are no visible lights, because everywhere the galaxies gleam far off. Still let’s move further away, always in a straight line, and for a hundred million years this time! We never collide with a barrier, and we never meet a wall, because there is neither limit nor edge nor end. The luminous points are always in front of us. We see some here and there, forming trails in space, but they are everywhere! It has now been a hundred million terrestrial years since we lost sight of our Galaxy, and we advance forever… forever… forever… The universe, it’s forever.
(13) By returning now slowly on Earth, we understand better what is eternity in this journey because, if it didn’t exist, with what could we limit the volume and the duration of the universe? With the word end? We just have to imagine limits to immediately question what forms them or on what is behind them, because the spirit can’t stop at limits. Also, everywhere we were, God was found. And it is with his eyes that we looked at the immensity during our journey which brought us next to myriads of living worlds, myriads of seas, filled with people, with vessels, and with prairies in which herds graze. But, over there, where was our reference to situate us in the immensity? What were our hour and our terrestrial year worth in these distant spaces? We were alone with the celestial Spirit, without any reference other than the luminous points extending into infinity, in all directions.
(14) Since we can’t make a volume of emptiness… it is certain that you didn’t travel within space that’s devoid of matter, but in the essence with which the spatial volume and the celestial bodies form. Also, listen to me! In the middle of all these wheels where you had gone, and which are composed of myriads of stars with their celestial bodies, did you really think that only one of these stars (our Sun) illuminated a living world? Did you have the feeling that all other stars shone vainly and pointlessly in the sky? When we shall demonstrate that it is the planets which make the stars shine, nobody will believe anymore that only one of these planets (ours) is inhabited among these myriads of wheels. This journey, which enlightens the mind on the depths of the universe, makes by itself become aware that the Sun isn’t the Star, or the Earth the Planet of the universe, but that they are only celestial bodies among the other similar celestial bodies which compose the Wheel.
Movement and displacement of the galaxies
(15) The rotation movement of our Galaxy can’t appear to us from the Earth. However, by being towards the edge (in the lower third) and far from the center of this big wheel, we move in space at an incredible speed! And the distances which we travel are terrifying! By moving this way, as on a circle, we are constantly moving away from certain galaxies and going towards to meet others. Those from which we are going away at a great pace seem to us necessarily more red than white, whereas those which we approach seem to us rather blue than white, because the speed of our movement makes their color vary for us. Now, this phenomenon of the change of color considerably deteriorates the defect of vision of the scholars. Because, by seeing the galaxies going away from them, while it is them who are going away from them, they can’t refrain from concluding that the universe is expanding... It is what they teach.
(16) By being near the extremity of a curved branch, we see more of the wheels from which we are moving away from than those we are going towards. And it is because of our own movement that the scholars see the universe as a gigantic ball constituted by galaxies which don’t stop moving away from one another, pushed by the breath of the initial explosion which would have formed them. But if it was so, wouldn’t the space where this expansion is occurring had to be unlimited? What has no limit is necessarily eternal. Only space would be eternal in their eyes? What is forming this space according to them, and how far do they imagine it to go? Is it a part of the universe or isn’t it a part of it? Shouldn’t they answer these questions?
(17) Besides, if we move away quickly and at will from their immense ball of galaxies (which is the totality of the universe for them), we would inevitably end up seeing it as big as a heap of stars, then as the Earth, then, further still, as a little point which would shrink even further until it dissapears from view... What that this tiny and strange universe in the middle of the unlimited space, and similar to a grain of sand which quickly disappears from view? Isn’t it a challenge for the reason?
(18) No, the universe isn’t a local and temporary little thing which would exist from an explosion of the nothing coming out of nowhere, and which would have accidentally occurred... It is God’s domain which has no limits of depth and of duration, because you saw that it exists everywhere and always, celestial travelers! You can’t thus have your own life, because you have the life of the Almighty by being one of its uncountable domains. As nobody has its own weight, no creature has its own life, its life being that of God which exists everywhere. This revelation will give you consistency and an evident reason to have hands and feet around a heart. Believe that it is so, because life exists wherever we are in the universe, and manifests itself on all the planets where there is some liquid water. Now, it will be shown that there are as many inhabited lands as there are stars! See then that we aren’t alone but numerous in the eternal immensity.
The proportions of the masses
(19) During our very long journey, we have noticed that, however big they are, the galaxies were in truth only tiny luminous points in the distance. And we also noticed that we were constantly in the center of a sphere bounded by these luminous points still perceptible to the eyes. Also, whether it is about a celestial body or about an entire galaxy, these bodies are really only corpuscles in the immensity. It is thus, because in the unlimited volume of the universe, the dimension of objects is only an affair of distance separating the objects from us. What is easily noticeable in our own Galaxy, where stars appear to us as tiny luminous points on the celestial vault.
(20) How would it then be possible to observe the planets of the stars from the Earth? We can’t do it, because even Jupiter appears to us as a luminous point in the sky, while it is at our doorstep. We see it because it is close to us. But if it was on the borders of the solar family, we would barely be able to see it, because a planet only weakly reflects the light of its star. Therefore it is inconceivable to be able to visually observe the planets of the stars, even the nearest ones. We cannot do it, especially when the disc of a planet is completely illuminated by its star (as the full moon) in relation to us, this planet is necessarily in the background of this star. Also such little light sent back by this tiny point is blended in the powerful light of the star which we observe. That is why, even with a powerful telescope, it is impossible to notice the planets of the stellar families others than the solar family.
The discretion of angels
(21) Knowing very well the matter and the force which compose together the universe, I say that, whatever its appearance is, the force is nothing more than matter in movement. There are not many forces, but the force. The waves are also a part of this unique force, because they are always vibrations caused in the matter, the diverse vibrations which echo here and there within this matter. It is also because space (the unlimited volume of the universe) is constituted of essence that the waves can spread and interpenetrate this space. Things which couldn’t be done with electrons, because, however small they may be, they are bodies which would collide against each another. No, the waves are only vibrations of the matter, and not emissions of electrons. We shall return to the waves later to explain them. And there you will grasp that all worlds of the sky are interconnected.
(22) If that’s how it is, you will say, why don’t the angels of the Wheel then make themselves known? What’s the reason they don’t show themselves, while we can communicate with them through waves today? My answer to that is that only the men of the Wheel wish to communicate with the other worlds, and not the angels who know the whole truth. Thus, in the Wheel, men’s worlds which are approximately in our time and still in the darkness, are necessarily at the opposite of us in the Galaxy and far away. There is therefore no possibility of communicating with them. Furthermore, wanting to communicate with other worlds is the expression of a great distress of the ignorant man which, still not knowing the truth, feels alone and abandoned in the immensity.
(23) Having nothing to say to them, the angels do not try to communicate with men. Isn’t it written that the Son would come to lead you in all the truth, at the appointed time? It is indeed necessary that it is him who leads you in it in this day, to separate the sons of darkness from the sons of light. He has to judge and put an end to the corrupt world; because the daily destruction of the Earth and his inhabitants has to stop to make place for the sanctuary and for God’s reign.
(24) When the time comes, the Shiloh (Christ) thus comes to put an end to all authority, all power and all domination on every Earth where God sends him. However, it isn’t the same man of flesh and bone that arrives in these worlds, but the spirit of truth, which is similar here or elsewhere in the universe and in the man whom God raises to serve him. The truth is unique. And if the words chosen to say it vary a little from one world to another, they do not express less the reality. And all hear.
(25) Thus, every world of the sky sees one day the Shiloh arriving with power, to stop with his finger the train of the world going into perdition. He does it through writing, by explaining patiently to men what reality is exactly. The only Son is unique only in the world where he rises. Because I have just said that, the Father has as many related sons as he has got living worlds in his universe. And his sons are alike in all respects, because God is the spirit which animates them, the unique spirit of truth.
(26) The worlds of the sky which are upstream of the Sun, have thus gone through this singular day of the coming of Christ. So it goes without saying that angels refrain from showing their presence, provided they don’t practice waves in the same manner as men does, and also because of what they could reveal, the Son reveals it. No, none of them would substitutes himself from the one to whom they owe their salvation, as will those who will pursue their life in the kingdom of God.
(27) If thus the angels of heaven would themselves explain the truth, they would annihilate necessarily the mission of the Son which consists in enlightening you, in separating you from goats, and in saving your lives. Do you understand that if they could communicate with the world with waves, they would do it necessarily with the scientists, the military, and the other leaders who dominate the peoples? In that case, these men by whom the end is coming would then be strengthened and taller than ever, whereas the poor people of humble heart would be forever their victims. There would be thus no purgatory. Then the princes of darkness would reign until they bring the extinction of all life. That’s why the angels of heaven don’t get involved in this world nor in any other worlds, because having experienced themselves the purgatory, they know that the Son will come at the appointed time to make the truth known to them and to save them from evil powers.
(28) Until his arrival, most men think they are alone in the immensity. Others imagine the universe populated with monsters waging war constantly, as do men of darkness from this world. But nobody thinks of the peaceful angels of both sexes knowing the truth, because they still don’t know what the reason to be of a celestial body is, or that saint men become angels at the appropriate moment. And such you were because, as princes, you had the teachers of lies and the organizers of the massive destruction of the world. How would you have been able then to know that the only monsters of the Wheel are those who destroy their Earth and its inhabitants? They are amongst you. They reign by holding you under their authority. And you know them!
(29) Since evil invalidates life and stops it, while good gives free rein and protects it, which of two must overcome the other? Are they the men who lead the world to its ruin who have to reign over the Earth or those who walk in the ways of the Almighty? If thus those who reign today could communicate with the worlds of the sky, or go from world to world as they hope to do, they would act on them as they act on you since the antiquity. What would happen then to an emerging world, as the one who will appear around Jupiter, if the leaders of peoples, the traffickers, the conquerors, the rich, the priests, the scientists and the military of our planet could access it? This distant world would become their victim as our world became one, and it couldn’t even make it to the days that we are into! That is why God doesn’t allow these individuals to communicate with other worlds, nor to travel in the Wheel. God rejects them, because they were only there to teach evil, and push it to its paroxysm and disappear forever.
(30) Know also how to see that what is practiced on the entire Earth today exists for the arrival of the Son and the change of world which follows. Nothing is missing in any domain. Everything is ready to serve me, and to serve He who sends me. That is why you will see me arriving everywhere in the world as fast as lightning, without anyone having the time to oppose to it and such as Jesus announces it.
*
(31) You, the saint angels, you are of the race of the Father in all the wheels, because God created you in his likeness. Don’t doubt what I say about your divine nature, because men who become angels are God together throughout the universe. And to better be aware of it, lay yourself bare, and go bathing again and several times in the intergalactic space, in the middle of the wheels filled with myriads of living worlds. You will come back with a clearer view and fairer thoughts on the dimensions of the universe, and with better feelings on the depths of He who gives you the being, the breath and the movement. | https://www.thebookoflife.eu/fulfilment-of-the-scriptures/the-worlds-of-the-universe.php |
> this? Has it worked? Is there a reason this wouldn't work?
I tried " -vf pullup,dejudder,idet,yadif=deint=interlaced,fps=60000/1001" and it seems to have worked. I added the fps=60000/1001 because my set top box can't deal with vfr or frame rates above 60fps. | https://lists.ffmpeg.org/pipermail/ffmpeg-user/2014-March/020714.html |
A longitudinal study of ABC transporter expression in canine multicentric lymphoma.
Canine lymphoma is typically treated with a doxorubicin-based multidrug chemotherapy protocol. Although this is often initially successful, tumour recurrence is common and frequently refractory to treatment. Failure to respond to chemotherapy is thought to represent drug resistance and has been associated with active efflux of cytostatic drugs by transporter proteins of the ATP-binding cassette (ABC) family, including P-glycoprotein (ABCB1), MRP1 (ABCC1) and BCRP (ABCG2). In this study, ABC transporter mRNA expression was assessed in 63 dogs diagnosed with multicentric lymphoma that were treated with a doxorubicin-based chemotherapy protocol. Expression of ABCB1, ABCB5, ABCB8, ABCC1, ABCC3, ABCC5 and ABCG2 mRNA was quantified in tumour samples (n = 107) obtained at the time of diagnosis, at first tumour relapse and when the tumour was no longer responsive to cytostatic drugs while receiving chemotherapy. Expression data were related to patient demographics, staging, treatment response and drug resistance (absent, intrinsic, acquired). ABC transporter expression was independent of sex, weight, age, stage or substage, but T cell lymphoma and hypercalcaemia were associated with increased ABCB5 and ABCC5 expression, and decreased ABCC1 mRNA expression. Drug resistance occurred in 35/63 (55.6%) dogs and was associated with increased ABCB1 mRNA expression in a subset of dogs with B cell lymphoma, and with increased ABCG2 and decreased ABCB8, ABCC1 and ABCC3 mRNA expression in T cell lymphomas. ABC transporter expression in the pre-treatment sample was not predictive of the length of the first disease-free period or overall survival. Glucocorticoids had no effect on ABC transporter mRNA expression. In conclusion, drug resistance in canine multicentric lymphoma is an important cause of treatment failure and is associated with upregulation of ABCB1 and ABCG2 mRNA.
| |
The invention concerns regulating the oxygen partial pressure of the gas mixture in the respiratory circuit of a diver, the diver being provided with a closed-circuit respiratory apparatus connected to a remotely disposed monitor chamber in which there is a pressure which is equal to or close to atmospheric pressure.
It is applied in particular to the case of divers connected to a submarine of the `diver-exiting` type, operating at a depth which can be down to 300 meters.
For divers to operate from a submarine, at great depth, it is necessary to use respiratory systems which have a low level of consumption as the amount of gas which can be stored on board submersible vessels is limited.
The respiratory systems of this type are `closed circuit` systems in which the oxygen partial pressure is measured for controlling the intake of oxygen. The devices which are known at the present day are of compact nature. Accordingly, the conditions in regard to pressure and humidity require the electronic measuring device to be of such a degree of complexity and to have components of such a quality that such systems are burdensome and require particular competence and supervision on the part of the diver. Now, it is certain that the first function of the diver is to work effectively. It is penalising him to require him to have a preoccupation with his respiratory system when he is in a situation in regard to pressure which affects his reflexes and his initiative and capacity for decision.
The present invention seeks to liberate the diver from all respiratory preoccupations by arranging for that supervision task to be performed by the personnel adjacent the diver in the atmospheric observation chamber of the submarine.
According to the invention, this is achieved by means of a process wherein:
(a) a sample of the gaseous mixture in the respiratory circuit of the diver is continuously taken off,
(b) said sample is passed into the monitoring chamber, the oxygen partial pressure of the sample is measured therein, and an electrical control signal is produced therein, if the measured pressure falls below a reference value, and
(c) the control signal is used for triggering an injection of a defined amount of an oxygen-rich gas mixture to the respiratory circuit of the diver.
By transferring the entire measuring and regulating arrangement into the atmospheric chamber, not only does the electronic apparatus operate under standard conditions in respect of pressure and humidity, but the breathing apparatus carried by the diver is reduced to a simple and compact device.
In a typical example, the sample of the gas mixture is taken off at a flow rate which is in the range of from 0.2 to 0.6 normal liter per minute.
In a preferred embodiment, the sample from the respiratory circuit is passed to the monitor chamber by a capillary conduit in such a way that the variation in the oxygen partial pressure at the outlet of the capillary reproduces the variation at that pressure at the inlet of the capillary, and therefore in the respiratory circuit, with a delay, the section of the capillary being so selected as to limit the duration of the transfer of the sample along the capillary.
In a preferred embodiment, the interval between two injections is used to prepare said defined amount.
In a preferred embodiment, said defined amount is prepared by filling a container with the oxygen-rich mixture at a given differential pressure.
In a preferred embodiment, said injection is produced by discharging, into a simple conduit leading to the respiratory circuit, a container which had been previously filled with oxygen-rich mixture, at a given differential pressure.
The invention is not limited to this use.
Thus, the breathing apparatus used by the diver to provide him with the gas mixture which he is to breathe, instead of being carried by the diver, may be a chamber or tank in which the diver is located and in which there is an atmosphere which is breathed by the diver.
In one use, the breathing apparatus is a chamber in which the diver is disposed and which contains a gas mixture which is breathed by the diver.
In one case, that chamber is a hyperbar chamber disposed in a location forming the monitor chamber and intended to simulate the diving depth.
In another case, the chamber is a working tank submerged at the working depth.
In another use, the diver operates in the water from a submerged chamber forming a diving chamber connected to a monitoring chamber disposed on board a surface vessel, and the gas mixture is prepared in the submerged chamber.
Various uses will be described hereinafter with reference to the Figures of the accompanying drawings in which:
FIG. 1 is a general diagrammatic view of the device,
FIG. 2 is a diagrammatic view of the respiratory circuit of the diver,
FIG. 3 is a diagrammatic view of a part of the device,
FIG. 4 is a diagrammatic view of a device according to the invention, in the case of a hyperbar chamber,
FIG. 5 is an alternative form of the device of FIG. 4,
FIG. 6 is a diagrammatic view of the device according to the invention, in the case of a working tank,
FIG. 7 is an alternative form of the device of FIG. 6, and
FIG. 8 is a diagrammatic view of a device according to the invention in the case of a diving bell.
FIG. 1 shows a `diver-exiting` submarine 1 which comprises a monitor chamber 2 in which the pressure is atmospheric and which contains the supervisory personnel, and a compartment 3 at the pressure of the bottom where the divers are diving. It is assumed that a diver 4 has left the compartment 3 and is operating in deep water, let us say at a depth of from 100 to 300 meters.
The diver 4 is provided with a breathing device (FIGS. 1 and 2) which forms a closed respiratory circuit comprising a mask or mouthpiece 5 connected by two pipes 6 and 7 to a flexible bag 8 with a capacity of 4 to 5 liters. Non- return valves 9 ensure that the gases circulate in the direction indicated by the arrows, in per se known manner. The circuit includes a cartridge 10 for fixing the carbon dioxide which is breathed out, and a valve 11 for the discharge of an overflow of gas, if appropriate.
According to the invention, this circuit (for example in the region of the bag 8) is connected to the submarine by an umbilical 12 which contains two conduits 13 and 14 communicating with the interior of the circuit.
The circuit 13 is a capillary conduit of very small section (of the order of 0.5 mm in diameter) which connects the circuit to a measuring apparatus 15 in the monitor chamber 2 of the submarine.
The apparatus 15 continuously measures the oxygen partial pressure of the gas taken off by the capillary 13 and, in dependence on that measurement, supplies an electrical signal which, by way of a connection 16, controls a three-way electro-pneumatic distributor means 17 in the diving chamber.
The electro-pneumatic distributor means 17 (FIGS. 1 and 3) has an output 17a to which the other conduit of the umbilical is connected, and comprises an inlet 17b connected to an external breathable gas tank 18 by way of a pressure-reducing valve 19 and another inlet 17c connected to a container 20 with a volume C, which is disposed for example in the diving compartment 3. The distributor means 17 communicates the container 20 either with the pressure reducing valve 19 or with the conduit 14.
In the rest condition of the distributor means, the container 20 is in communication with the pressure reducing valve and is charged to a stabilised differential pressure &Dgr;p. In operation of the distributor means, it is communicated with the conduit 14 into which it is discharged, supplying the diver with an amount C &Dgr;p of breathable gas. The value of &Dgr;p is pre-regulated and defined in dependence on the oxygen content of the gas mixture used.
Thus, in each discharge operation, the amount of oxygen supplied to the diver is perfectly defined. The respiratory circuit of the diver is closed between two injections.
The measuring device 15 is controlled to supply an electrical signal which controls the distributor means in order to discharge the container 20 into the conduit 14 as soon as the oxygen partial pressure detected is lower than a minimum reference value.
For example, according to the invention, the conditions are so set that, in each discharge, the oxygen partial pressure of the gas mixture which is breathed by the diver is raised by a constant value of 200 MB. The circuit of the diver is closed between two injections and the pressure fluctuates between the minimum reference value (for example 300 mb) and that value when increased by 200 mb (that is to say, 500 mb). The frequency of the injection operations depends on the work being done. At rest, it is about 1 injection every 2. 5 minutes.
In order to eliminate any danger of hyperoxia, the oxygen content of the injected rich mixture is so selected that, at the working depth, the oxygen partial pressure of the injected mixture is close to 1000 mb. If required, the diver can then use that mixture freely by means of a direct feed connection 21 which connects the respiratory circuit of the diver to the pressure reducing valve 19, by means of a manual control 22; in particular, this manually controllable feed makes it possible if required to adjust the current respiratory volume.
The `oxygen-rich` gas mixture is a mixture which is richer than the reference value. In the limit condition, the mixture may be pure oxygen.
The device is most suitable for diving depths of from 20 to 300 meters, the capillary providing an insufficient flow at a depth of less than 20 meters and an excessive response time at a depth of more than 300 meters.
In the embodiments shown in FIGS. 4 to 7, the device comprises:
a chamber D with diver therein;
a capillary gas conduit 13 between the chamber D and a device 15 for measuring the oxygen partial pressure, such device being disposed in a monitor chamber in which the pressure is close to atmospheric pressure, the device 15 being capable of measuring the oxygen partial pressure of the gas mixture in the chamber in which the diver is located and supplying a control signal when that pressure falls below a reference value;
a gas supply conduit 14 which opens into the chamber containing the diver; and
means G for preparing a defined amount of an oxygen-rich gas and injecting said amount into said supply conduit under the control of the control signal.
The means G represent the assembly of the means 17 (distributor means), 19 (pressure reducing valve) and 20 (container) of the device of FIG. 1, or equivalent means.
In the embodiment of FIG. 4, the chamber D is a hyperbar chamber and the means G are outside the chamber, the exterior of the chamber forming the monitor chamber.
Preferably, the pressure reducing valve 19 is pilot-controlled by the pressure within the chamber D, the pilot-controlled connection being shown by the broken line connection shown in the drawing.
In the embodiment of FIG. 5, the means G are disposed within the hyperbar chamber D (decompression tank or saturation chamber).
In the embodiment of FIGS. 6 and 7, the chamber D forms a working tank which is positioned directly on the site of operations and the diver operates therein in a dry condition, without a breathing mask, so that he directly breathes the atmosphere in the chamber.
In the embodiment of FIG. 6, the pneumatic system is disposed in a monitor chamber forming part of a surface vessel N. The chamber is connected to the working tank by the capillary conduit 13 which connects to the measuring device 15 and by the supply conduit 14 which connects the interior of the working tank to the pneumatic system G.
In the embodiment of FIG. 7, the pneumatic system G is disposed in the tank itself and the monitor chamber only contains the measuring device 15.
In the embodiment of FIG. 8, the device comprises:
a closed circuit respiratory device 8 worn by the diver;
a monitor chamber which is disposed on board a surface vessel N and in which there is a pressure close to atmospheric pressure;
a capillary gas conduit 13 between said circuit and a device 15 for measuring the oxygen partial pressure, which is disposed in the monitor chamber and which is capable of measuring the oxygen partial pressure of the gas mixture in the circuit and supplying a control signal when that pressure falls below a reference value; and
a submerged diver chamber P provided with means G for preparing a defined amount of an oxygen-rich gas and injecting said amount of gas into a gas supply conduit 12 connecting the diver chamber to the respiratory circuit of the diver, under the control of the control signal. | |
Please read the conditions and guidelines required for donations and loans in the Collections Policy below. If you still wish to donate or loan artifacts, please describe them on the Certificate of Gift Form or the Agreement to Loan Property Form, whichever is applicable. Please print out the completed form and bring it to the museum with your items.
Vintage Wings and Wheels Museum Collections Policy
The Poplar Grove Vintage Wings & Wheels Museum was established to preserve history and educate the public about the significant contributions made to our country by winged and wheeled vehicles. The artifacts the Poplar Grove Vintage Wings & Wheels Museum accepts must be related to Early Transportation History. To achieve these goals, the scope of our collection shall be limited to items that enhance our mission statement and that Poplar Grove Wings & Wheels can ethically and physically provide care for. Specifically, we are looking for:
- Pre-1957 artifacts.
- Historically significant materials that fit within the scope of the museum’s collection and further its mission.
- Actual historic items, not merely a copy or reproduction of an historic item.
- Items in good condition and ones that the museum can adequately care for.
*Items that are in very poor condition, are broken beyond repair, moldy, have been/are infested with insects or other vermin will not be accepted. *
This includes but is not restricted to artifacts related to:
- Aviation, private and military
- Vehicles such as: automobiles, bicycles, tractors, wagons, carriages, and more.
- Related photographs and print media.
- Related ephemera.
The Poplar Grove Vintage Wings & Wheels Museum gladly accepts item donations for our collection that augment our current collection. Items in the museum’s collections are carefully stored and preserved and are used to enhance museum exhibits and educational programs. Prospective donors should understand that we may have to refuse certain gifts or loans for any of the following reasons.
- Duplication of currently owned material
- Lack of space
- Lack of technical equipment or financial resources to properly care for or process material offered
- Inappropriateness to the collection, such as outside of our scope and mission purpose
* Any Materials’ that fall outside these proscribe limits will be reviewed by the Collections Committee on a case-by-case basis. *
Acquisitions to The Poplar Grove Vintage Wings & Wheels Museum collections by purchase, loan, gift, bequest or other means shall accord with the following rules:
- The owner must have clear title and must sign a deed of gift transferring title to the Poplar Grove Vintage Wings & Wheels Museum. In the case of a bequest, the donor must also have had clear title.
- A transfer of ownership file containing gift agreements and other proofs of the Poplar Grove Vintage Wings & Wheels Museum’s legal ownership of acquisitions shall be maintained.
- Vintage Wings & Wheels Museum does not do monetary appraisals. (See U.S. Internal Revenue Service regulations.)
- The Poplar Grove Vintage Wings & Wheels Museum must be capable of housing and caring for the proposed acquisition according to accepted professional standards.
- Proposed acquisitions shall be free of donor-imposed restrictions unless such restrictions are agreed to by the Museum staff.
- The Museum staff reserves the right to refuse gifts or portions of donations, and the right to dispose of or return to the donor, items inappropriate to the collection.
- Acquisitions approved by the Museum staff shall be promptly accessioned upon receipt and processed.
- Donors and prospective donors, whenever deemed appropriate, should be asked by the Museum staff whether they would be willing to provide funds for the full or partial cost of accessioning and subsequent maintenance of materials gifted to the Museum. Willingness or unwillingness to provide such funds should usually not be a determining factor to accept or reject a gift for accessioning.
Two types of donations:
v A gift is a donation that becomes The Poplar Grove Vintage Wings & Wheels Museum full property, including title, property rights, trademarks and copyrights. The Poplar Grove Vintage Wings & Wheels Museum prefers to be the sole owner of any material accepted. This allows The Poplar Grove Vintage Wings & Wheels Museum to retain absolute authority to use the item(s) in any way we deem appropriate. Thus, restriction on the handling or use of donations which impose an unsupportable technical, financial or professional burden on The Poplar Grove Vintage Wings & Wheels Museum cannot be accepted. Reasonable restrictions on a gift may be accepted by the Museum staff. The Vintage Wings & Wheels Museum reserves the right to dispose of material in deteriorating condition or under changed museum circumstances in compliance with accepted museum procedure and professional responsibility.
v A loan is a donation that continues to belong legally to the lender but is temporarily housed at The Poplar Grove Vintage Wings & Wheels Museum and may be used for exhibits and displays at the Museum’s discretion. The lender has the right to take the loaned item back from the Museum with sufficient notice. At the end of the loan period, if the lender decides not to renew, the loan period will end, and the item must be removed. If the Museum decides not to renew the loan period, it will contact the owner and ask that the item be removed, and the loan period will end. At any time, the owner may gift the item to the Museum and will fill out a Deed of Gift. All items on loan will be reviewed by November 1, each year. The Museum requires valuable property to be insured by the owner for the life of the loan.
*These policies do not include the Library Collection. | https://www.wingsandwheelsmuseum.org/artifact-donation/ |
Celestial Mechanics and Astrodynamics
Research at CSR applies the principles of physics to the precise determination, prediction, and optimization of trajectories in space. Activities include the characterization of the motions of the Earth, Moon, and other celestial objects, as well as of rockets and artificial satellites, both terrestrial and interplanetary. Applications include precision determination of spacecraft orbit and attitude dynamics, mission trajectory design and operations from launch and navigation to re-entry and landing.
Data Exploration: Models, Algorithms and Error Analysis
CSR research centers on the use of large global-scale data sets to solve complex and computationally challenging problems and to address open-ended questions demanding advanced and innovative approaches. Using some of the world’s most powerful high performance computers, CSR has developed and employed innovations in advanced high-performance modeling, estimation techniques, pioneering statistical methods and error analyses.
Space Geodesy
CSR is a globally-recognized leader in the measurement and representation of the Earth and other celestial objects, including gravitational fields in a three-dimensional time-varying space for positioning, navigation, and observing geodynamical phenomena. Resident capabilities encompass coordinate systems, control networks and control techniques for studying oceanography, hydrology, and geophysics.
Mission & Data Architecture, Design and Simulation
CSR routinely contributes to space mission design initiatives its expertise in satellite navigation, attitude and scientific payload hardware configuration, simulation and testing, space and ground system implementation, data products design, hardware and software networking configuration, and delivery and verification of onboard navigation software and models.
Radionavigation
CSR has conducted research with the national GPS, and more generally the global GNSS, since the inception of both. Current research in this domain includes precision and opportunistic navigation, navigation system protection, the design of software-defined receivers and other technology innovations, and the study of the ionosphere and neutral atmosphere.
Remote Sensing
Since the 1970s, CSR has been deeply involved in the science and art of identifying, observing, and measuring planetary topography from afar. CSR’s altimetry expertise began with early satellite radar altimeters and continues through more recent American and European missions. CSR is also at the forefront of both airborne and space-borne laser altimetry (LIDAR) systems and system design.
Satellite Technology
Mission planning and associated hardware development has long been an important research thrust at CSR. CSR’s expertise includes spacecraft planning and design, launch, GN&C (guidance, navigation, and control), operations, and decommissioning/closeout. In 2002, with the goal of inspiring students through hands-on participation, the UT Satellite Design Laboratory (SDL) was founded to provide students with an end-to-end experience of space technology and mission operations.
Scientific Interpretation and Analysis
Results of CSR’s research provide solutions to questions associated with many disciplines and cross-disciplinary fields including fisheries and aquaculture, agriculture, weather forecasting, environmental impacts of oil spills, oil exploration and drilling operations, mapping ocean circulation, sea-level rise and ice sheet decline, and significantly improved models of the Earth’s gravity field.
CSR continually strives to expand and deepen its contributions towards addressing complex scientific challenges, such as quantifying sea level rise and ice cap melt, monitoring global water storage processes, developing innovative GPS/location-based applications, fabricating small satellites, and pushing each new supercomputer to its limits. | http://www.csr.utexas.edu/research/ |
How long does your baby sleep?? That's one of recent thoughts lately as I'm struggling to let Leia sleep on her usual routine prior to reaching 1 year old when she takes two naps. Now, she (almost) doesn't really want to sleep in the morning (after breakfast and play), unless she becomes easily tired depending on the activities in the morning. So I found out that babies at this stage might drop one of their naps and just sleep longer after lunch. There goes my worry then, just because what Leia does is normal.
Maybe babies still are adjusting, or maybe they still go through stages in which they're still finding their comfort zone. Babies are different, some sleep through the night while others don't like bedtime. But most kids share a a common tendency to feel sleepy at one time also. But Moms should make be able to provide a sound sleeping environment so that our babies will know right away when it's time for bed.
Ideally, a toddler should sleep 10-14 hours a day. It should include a nap of at least 2 hours or maybe 2- 1 hour naps. Recently Leia mostly does the first, after lunch or we call it siesta, she always sleeps for 2 hours or even more. Sometimes I fear it that what if there comes a time that she'll sleep less or take no naps. But I'm always thinking that ever since I could remember, I never had any problems with Leia's sleeping routine. Ever since she's a newborn she always sleeps soundly and rarely wakes up to feed during the night. Being a CS Mom, I did not wake her up in the middle of a sound sleep since I feel like sleeping is equally important with eating.
Based on research and as observed with Leia, sleep time should more or less be controlled or made a routine if you want to achieve healthy sleeping habits with your baby. In the morning, for example, if your baby was able to sleep after breakfast or play, sleep should less than 1 hour only. Then after lunch, for her second sleep time, don't let your baby sleep for more than 2 hours so you won't have hard time putting her back to bed for the night.
I believe that it helps babies sleep well if we let them know and feel that it's time to go to bed already. What do I do to create a sleeping environment for Leia? Here are some of the practices we do at home to let her sleep better.
1. I make sure she's well-fed, she's eaten enough along with milk.
2. She's clean, not sweaty. Calm and not hyped.
3. I read her a bedtime story every time before bedtime. Her favorite, Goodnight Moon.
4. Then, I always let her sleep in her crib. The crib should be intended only for sleep and not for play.
5. There's lullaby playing in the background or me singing until she sleeps soundly
6. She's with her hotdog pillows or her teddy bear.
7. Very important for Leia to sleep longer is be in an air-con room that's preferably without any loud noise.
8. During the night, we only have a small dimmed light so generally the room is dark.
9. It's best that I stay in the room when she sleeps because I can easily put her back to sleep in case she wakes unintentionally
10. Important to never leave you baby unattended. Have someone look over in case you need to slip out of the room for a short while.
So in case, some of these don't work out for you, give it time to let your baby adjust and work with the routine. Generally, these will help you cope with healthy sleep patterns with your baby. Nevertheless, do whatever it takes as long as you know that it's for the best for your baby. Sleep is very important because babies have a long way to go in their developmental stage. And this too will help them be healthy and active as they grow bigger and older. | https://www.leiasmom.com/post/2017/09/30/baby-sleep-101 |
Become A Wavemaker
Report | August, 2004
The appearance this summer of several cetaceans stranded on the coasts of the Canary Islands and the Azores while naval manoeuvres were being carried out has reopened the debate on the impact on cetaceans from the use of sonar and other acoustic pollution arising from these exercises.
This is not the first time this has happened in the Canary Islands, nor is this the only region in the world where the death of cetaceans has coincided with warship manoeuvres.
Despite the fact that the navies involved have repeatedly tried to deny their responsibility in these events, the fact is that both NATO and the US Navy have been aware of the cause of these deaths for years.
What is LFAS?
LFAS, or SURTASS LFAS, is the acronym for the high-precision SONAR system known as Surveillance Towed Array Sonar System, Low Frequency Active Sonar.
It is based on the use of high intensity sound waves (over 200 dB) at low frequency (between 450 and 700 Hz) that can travel great distances underwater and detect objects hundreds of kilometres away. Dozens of them are emitted in a matter of seconds (up to 250 within 4-5 seconds) and they hit objects and rebound to a receiver that interprets them and allows the object in question to be visualised. Sonar can also be used for a minute or more at a time at intervals of 10-15 minutes. This sound transmitter is suspended from the ship at a depth of around 50 metres.
However, it is known that NATO is continuing to experiment with systems at even lower frequency (50-150 Hz) and at a range of 230 dB, which would allow them greater reach and precision. Mid-frequency sonar is also being used, sometimes in combination with LFAS, with similar harmful effects. For this reason, both mid and low frequency sonar have been appointed as the cause for cetacean strandings. But the potential impact of LFAS is higher due to its range, and that lower frequencies can interact with whales’ sounds. Sound travels 4.5 times faster in water than in air and the lower the frequency (Hz) the further it can travel (hundreds of kilometres). In addition, the intensity (dB) is more consistent. Frequencies below 1 kHz lose barely 0.04 dB per kilometre.
| |
Q:
How do I ensure an animation is complete before initializing another function or animation?
so here, I have a sequence of animations using Raphael:
fade in curve
fade in ball 1
animate ball 1
fade in ball 2
animate ball 2
however, with my code, steps 4-5 are initiated WHILE the steps 2-3 are still animating. How do I ensure steps 4 and 5 are initiated after the animations of 1-3 are complete? I've tried using setTimeout on my second function (ball2), but no luck.
View on JSFiddle or here:
Raphael("bounce", 640, 480, function () {
var r = this,
p = r.path("M0,77.255c0,0,269.393,37.431,412.96,247.653 c0,0,95.883-149.719,226.632-153.309").attr({stroke: "#666", opacity: 0, "stroke-width": 1}),
len = p.getTotalLength(),
e = r.circle(0, 0, 7).attr({stroke: "none", fill: "#000", opacity:0}).onAnimation(function () {
var t = this.attr("transform");
});
f = r.circle(0, 0, 7).attr({stroke: "none", fill: "#000",opacity:0}).onAnimation(function () {
var t = this.attr("transform");
});
r.customAttributes.along = function (v) {
var point = p.getPointAtLength(v * len);
return {
transform: "t" + [point.x, point.y] + "r" + point.alpha
};
};
e.attr({along: 0});
f.attr({along: 0});
var rotateAlongThePath = true;
function fadecurve(ca,ba,aa,ab){
ca.animate({opacity:1},500);
setTimeout(function(){fadeball(ba,aa,ab);
},1000);
}
function fadeball(ba,aa,ab) {
ba.animate({opacity:1},400);
setTimeout(function(){run(ba, aa,ab);
},1000);
}
function run(ba,aa,ab) {
ba.animate({along: aa}, ab, ">", function () {
ba.attr({along: aa});
});
}
function startbounce() {
fadecurve(p,e,.9,400),
setTimeout(function(){fadeball(f,.8,400);
},1000);
}
startbounce();
});
A:
According to Raphael's documentation, the animate method takes a callback method as it's fourth argument. That method could be used to initiate the next animation in your sequence (or after the third animation).
function fadecurve(ca,ba,aa,ab){
ca.animate({opacity:1},500,,function(){fadeball(ba,aa,ab);});
}
For example.
| |
Every year many cockatiels’ eggs fail to hatch. I’m often asked “Why are my cockatiel eggs not hatching?” and in many cases, the problem is that the birds are young and inexperienced and simply don’t know what to do, or that they don’t have a suitable place for nesting. There are ways you can help your birds lay more fertile eggs, whether they’re currently well-nesting or not.
If you’ve got a cockatiel in need of hatching assistance, this post is for you. I’ll outline what causes infertile eggs, and offer tips on how to make sure your bird gets some fertile ones. It’s time to start thinking about the breeding season – these tips could be just what your cocky needs.
What Causes Infertile Cockatiel Eggs?
There are several reasons why cockatiel eggs are infertile. As the egg itself is an unfertilized egg or the bird’s age is not yet mature enough. Even the lack of material to make a nest can affect the infertility of the eggs themselves.
Here are some explanations for why your cockatiel eggs are infertile:
Unfertilized Eggs
A few cockatiels lay infertile eggs even when they are fertile. Their eggs are simply not fertilized – the sperm never made it to where it was needed. Usually, these eggs are found in nests that aren’t a strong enough environment for raising young birds – often due to new owners who don’t yet have the birds well-established.
In fact, it is not uncommon for cockatiels to lay eggs in the feed or at the bottom of their cages. This usually happens when you only keep a female without a male partner.
Age
Cockatiels are not mature enough to breed until about one year of age. Sometimes young birds will lay, but the eggs will be infertile. Fortunately, once your birds have reached breeding maturity, you can help them get fertile eggs more often.
Lack of Genetic Diversity
In order for an egg to be fertile, it must receive genes from both parents. If a pair of birds are closely related (e.g., bred from brother and sister), they may have an easier time than unrelated pairs, but some still won’t produce fertile eggs at all. If you’ve got birds of the same species, this is less likely to be a problem because their genes are more diverse.
Health
The health factor of the bird itself also greatly affects the fertility of the eggs.
As you know, cockatiel eggs need several nutrients which are crucial in their development. Therefore, you are obliged to provide healthy food when the breeding season arrives. Providing egg food can increase the chances of getting fertile eggs.
You may also read the following articles:
Lack of Nesting Material
Without adequate nesting material, birds may not be able to create a proper nest for laying eggs or raising chicks. Cockatiels typically need tree branches and shredded paper or cotton to build a suitable nest – if they don’t have what they need, they will sometimes lay eggs elsewhere, like on the bottom of the cage or in a corner where the cage touches the wall.
Too Young of an Egg
This can happen with inexperienced breeding pairs, in which case it’s not usually a problem. However, an egg that is too young may not have enough time to develop into a chick before it expires.
How to Get More Fertile Eggs
Don’t assume your birds have poor nesting habits because they’ve never laid fertile eggs before! Your birds may simply need to learn how to properly nest first.
Try the following to help your birds lay fertile eggs:
Sufficient Cage Size
Make sure there are at least three feet of cage space for them. This allows room for nesting material, some movement from the birds themselves, and plenty of space in case you want to put your bird in a different cage during the breeding season.
This allows room for nesting material, some movement from the birds themselves, and plenty of space in case you want to put your bird in a different cage during the breeding season.
Nest Box
Put more than one nest box in the cage. Cockatiels can be territorial and may prefer to move outside their nest boxes instead of using them for privacy. Also, you can sometimes see a female cockatiel’s eggs before they are fertilized by detecting the outline of a nest – she will often sit on top of her eggs for this reason.
If you have only one box, your bird may be less likely to use it. A good rule of thumb is to make sure there are at least two boxes per cockatiel on average during mating season.
Cockatiels can be territorial and may prefer to move outside their nest boxes instead of using them for privacy. Clean out old nest material as needed and replace it with fresh tree branches or paper towels.
If you have a “nest box” cage that uses plastic mesh, check the box on a regular basis to make sure it has no holes in the bottom.
Give Healthy Foods
Make sure your birds don’t have any fatty foods, like avocado or cheese. These can slow down their reproductive systems and keep them from breeding.
If you have some of your bird’s food in the cage, consider using a different type. Some foods contain oils that can dissolve feathers and make them fall off. By removing these oils, you can help prevent feather plucking and other stress-induced behavior problems.
Put your bird on a high-quality pelleted diet for breeding season to ensure good health and plenty of nutrients for laying eggs.
Keep it Clean
Clean any moss from the cage on a regular basis, as it can get moldy and mold can put stress on birds’ reproductive systems.
A New Nest Box is Better
If your bird has an older nest box, give it a new one to replace that one! Birds are known for breeding in anything from cardboard boxes to toy jars – so if your bird is used to nesting in that particular box, it’s probably better off with a new one.
Age is Matter
Never breed young birds (i.e., those under one year old). There is often too much of a chance that they won’t be healthy enough to handle breeding, or will otherwise not be mature enough to do so successfully.
If your bird has never laid fertile eggs before, give it at least two or three months before trying to breed it. This is a good amount of time for it to learn how to nest properly, as well as time for you to make sure the other birds in your home are able to lay eggs successfully.
Conclusion
As we have seen, there are several reasons why cockatiel eggs do not hatch. The most influencing factor is the readiness of the bird itself.
But don’t worry, the solutions above are an easy one to implement. And also patience is needed in this matter.
If your bird continues to lay infertile eggs even after trying the above tips, do consider that it may also be a sign of an underlying health or stress issue. If the problem is the health of your cockatiel, then the best solution is to make a visit to the vet.
You may also read the following articles: | https://parrotbrainery.com/cockatiel-eggs-not-hatching/ |
(Download PDF/Epub) Methods In Clinical Phonetics By Martin J. Ball, Orla M. Lowry(auth.)
This book is written for the beginning student of communication disorders with a basic understanding of phonetics, or the practising speech-language therapist whose phonetic training may need updating. It introduces the reader to the main areas of phonetics, and the main methods through which the phonetician reduces speech data to a permanent record.
The book, then, illustrates the three main approaches to the investigation of spoken language; articulatory, acoustic, and auditory. Further, it describes how impressionistic phonetic transcription through symbolisation differs from instrumental phonetic techniques. For each of these areas of discussion, chapters are provided that examine the general phonetic aspects, followed by chapters that illustrate their application to clinical data.
The authors are both phoneticians with experience of investigating both normal and disordered speech through both impressionistic and instrumental means, and this is the first book in this market that describes a whole range of data reduction techniques and illustrates them with data relevant to the student and practitioner of communication disorders.Content:
Chapter 1 What is Clinical Phonetics? (pages 1–9):
Chapter 2 Transcribing Phonetic Data (pages 10–24):
Chapter 3 Transcribing Disordered Speech (pages 25–40):
Chapter 4 Articulatory Instrumentation (pages 41–48):
Chapter 5 Articulatory Analysis of Disordered Speech (pages 49–60):
Chapter 6 Acoustic Instrumentation (pages 61–72):
Chapter 7 Acoustic Analysis of Disordered Speech (pages 73–87):
Chapter 8 Auditory and Perceptual Instrumentation (pages 88–98):
Chapter 9 Auditory and Perceptual Analysis of Disordered Speech (pages 99–109):
Chapter 10 The Future of Clinical Phonetics (pages 110–120): | https://1billionbooks.com/methods-in-clinical-phonetics-by-martin-j-ball-orla-m-lowryauth-download-ebook/ |
The purpose of this lesson is for students to represent multiplication situations with arrays and multiplication expressions.
Lesson Narrative
In a previous lesson, students arranged objects into arrays and described the arrays in terms of equal groups. In this lesson, students write expressions to represent arrays to further connect arrays and multiplication (MP2).
As students connect arrays to expressions, they may write \(3\times5\) or \(5\times3\) to represent 3 rows of 5 chairs. This is fine as long as students can correctly describe where the “3 rows of 5 chairs” are in their array or expression. Keep collecting ideas that arise about commutativity.
- Representation
- MLR2
Learning Goals
Teacher Facing
- Represent multiplication situations with arrays and multiplication expressions.
Student Facing
- Let’s represent situations with arrays and expressions.
Required Materials
Materials to Gather
Required Preparation
Activity 1:
- Each group of 2 will need 20 connecting cubes or counters.
CCSS Standards
Addressing
Lesson Timeline
|Warm-up||10 min|
|Activity 1||20 min|
|Activity 2||15 min|
|Lesson Synthesis||10 min|
|Cool-down||5 min|
Teacher Reflection Questions
In an upcoming lesson, students will learn about the commutative property of multiplication. What do you notice in their work from today’s lesson that you might leverage in that future lesson? | https://im.kendallhunt.com/k5/teachers/grade-3/unit-1/lesson-18/preparation.html |
In spite of the fact that the extraordinary progress of experimental techniques make us able to manipulate at will systems made of any small and well defined number of atoms, electrons and photons - making therefore possible the actual performance of the gedankenexperimente that Einstein and Bohr had imagined to support their opposite views on the physical properties of the wavelike/particlelike objects (quantons) of the quantum world - it does not seem that, after more than eighty years, a unanimous consensus has been reached in the physicist's community on how to understand their "strange" properties.
Unfortunately, we cannot know whether Feynman would still insist in maintaining his famous sentence "It is fair to say that nobody understands quantum mechanics". We can only discuss if, almost thirty years after his death, some progress towards this goal has been made. I believe that this is the case. I will show in fact that, by following the suggestions of Feynman himself, some clarification of the old puzzles can be achieved. This chapter therefore by no means is intended to provide an impartial review of the present status of the question but is focused on the exposure of the results of more than twenty years of research of my group in Rome, which in my opinion provide a possible way of connecting together at the same time the random nature of the events at the atomic level of reality and the completeness of their probabilistic representation by the principles of Quantum Mechanics.
1.2. The two slits experiment
In order to introduce the reader to the issues at stake I will briefly recall the essence of the debate between Bohr and Einstein which took place after the Fifth Solvay Conference (1927) where for the first time the different independent formulations of the new theory were presented by Heisenberg, Dirac, Born and Schrödinger, together with their common interpretation by Bohr - the socalled “Copenhagen interpretation” of Quantum Mechanics - which won since then a practically unanimous acceptance by the community.
This acceptance remained unquestioned for thirty years until when the books by Max Jammer (Jammer a1966, b1974) presented again to the new generation of physicists the ambiguities which still remained unsolved, and stimulated a renewed interest on those conceptual foundations of the theory which had been set aside under the impact of the the extraordinary experimental and theoretical boom of physics triggered at the end of World War 2 by the opening of the Nuclear Era.
The central issue of the debate, according to Jammer’s reconstruction (Jammer b1974 p.127), was “whether the existing quantum mechanical description of microphysical phenomena should and could be carried further to provide a more detailed account, as Einstein suggested, or whether it already exhausted all possibilities of accounting for observable phenomena, as Bohr maintained. To decide on this issue, Bohr and Einstein agreed on the necessity of reexamining more closely those thought-experiments by which Heisenberg vindicated the indeterminacy relations and by which Bohr illustrated the mutual exclusion of simultaneous space-time and causal descriptions.”
The thought experiment which both agreed to discuss was the diffraction of a beam of particles of momentum p impinging perpendicularly on a screen D with two slits S1 and S2 at a distace d from each other. Each particle, which passes through, falls, deviating at random from its initial direction, on a photographic plate P located after the screen. When a sufficiently high number of particles has been detected, a distribution of diffraction fringes typical of a wave with a central maximum and adjacent minima and less pronounced maxima appears. Each particle is detected locally, but seems to propagate as a wave.
Its wavelike nature is expressed by the Bragg’s relation connecting the wavelength λ of the wave in terms of the the distance d between the slits and the angle φ subtended by the central diffraction maximum (λ = φd). On the other side its particlelike nature is expressed by its momentum p which is connected to the the wavelength by the de Broglie’s relation (p = h/λ).
Since it is not possible to detect through which slit the particle is passed, its position x on D is uncertain by ∆x = d. For the same reason, the momentum acquired by the particle in deviating from its initial direction normal to D is uncertain by ∆p = φp
From these relations the Heisenberg uncertainty relation
follows. Incidentally, the same phenomenon occurs with only one slit, with d now indicating the slit’s width.
For Bohr eq. (1) holds for each individual particle. The particle’s position x and its momentum p are, in his words, “complementary” variables. They cannot have simultaneously well defined sharp values. In the interaction with the classical instrument made of the screen D and the photographic plate P, each particle of the beam acquires a blunt value x affected by an uncertainty ∆x and a blunt value p affected by an uncertainty ∆p. The product of the uncertainties however, can never be less than the limit set by (1). Initially, before impinging on the instrument, each particle was in a state with a well defined sharp value of the momentum and a totally non localized position in space. At the end, after having been trapped in the photographic plate, each particle has acquired a well defined sharp value of its position in space, and has lost a well defined value of the momentum. The essence of the argument is that only by interacting with a suitable classical object one side of the quantum world acquires a real existence, at the expense of the complementary side becoming unseizable.
For Einstein instead Quantum Mechanics is only a statistical theory which does not fully describe reality as it is. The uncertainties, according to him, reflect only our uncomplete knowledge. He postulates the existence of “hidden variables” of still unknown nature, and concentrates his efforts on proving that Quantum Mechnics is “incomplete”. In fact - he argues - if D is not fixed but is left free to move, one could identify the slit through which the particle has passed by measuring the recoil of the screen produced by the momentum exchange with the particle deviated from its straight path. Both the position and the momentum of the particle could in this way be measured, violating the Heisenberg limit.
This does not work, however - replicates Bohr ( Bohr 1948) - because the detection of “which slit” changes the diffraction pattern. In fact, he argues, if, by detecting the recoil of the screen one determines through which slit the particle has passed, the position in space of D becomes delocalized by a quantity ε in such a way that the resulting maxima and minima of the possible two-slit diffraction patterns superimpose and cancel each other. The original diffraction pattern with D fixed becomes the diffraction pattern of the single slit through which the particle is passed. ∆x is reduced to the width of the slit and the uncertainty ∆p is correspondingly increased. Heisenberg’s relation for the particle still holds.
“It is not relevant - Bohr wrote many years later (Bohr 1958a) in a report of his debate with Einstein - that experiments involving an accurate control of the momentum or energy transfer from atomic particles to heavy bodies like diaphragms and shutters would be very difficult to perform, if practicable at all. It is only decisive that, in contrast to the proper measuring instruments, these bodies, together with the particles, would, in such a case constitute the system to which the quantum mechanical formalism has to be applied.”
On the other hand, Bohr insists to stress the classical nature of the instrument (Bohr 1958b): "The entire formalism is to be considered as a tool for deriving predictions of definite statistical character, as regards information obtainable under experimental conditions described in classical terms.[..] The argument is simply that by the word “experiment” we refer to a situation where we can tell others what we have learned, and that, therefore, the account of the experimental arrangement and the results of the observations must be expressed in unambiguous language with suitable application of the terminology of classical physics."
It is therefore clear that for Bohr the proper measuring instruments on the one side must be treated as classical objects, but on the other one that the parts of the apparatus used for the determination of the localization in space time of particles and the energy-momentum transfer between particle and apparatus must be submitted to the quantum limitations. We will come back in a moment to this question in order to prove that this ambiguity can be understood in the framework of an interpretation of Quantum Mechnics in which both Einstein’s purpose of saving the objectivity of the properties of macroscopic objects and Bohr’s denial of the possibility of attributing to the objects at the atomic level independent properties, are recognized.
1.3. The EPR paradox
The second phase of the debate sees a change in Einstein’s strategy of proving that the description of reality given by Quantum Mechanics is incomplete. This phase is based on the formulation of the EPR (Einstein, Podolski, Rosen) paradox (Einstein et al 1935).. I will briefly sketch its main argument, even if it is not essential for the further development of the argument of this Chapter.
This is how the authors formulate the basic assumption of their argument: "If, without in any way disturbing a system, we can predict with certainty the value of a physical quantity, then there exists an element of physical reality corresponding to this physical quantity."
Consider a system of two particles in a state in which the relative distance x1 - x2 = a and their total momentum p1+p2 = p are fixed. This is possible because these quantities are not complementary. Then EPR argue as follows. By measuring the position x1 of the first particle it is possible, without interfering directly with the second particle, determine its position x2= a+x1. This means that, according the initial definition, that x2 is an element of reality . However, we might have chosen to measure, instead of x1 the momentum p1 of the first particle. This measurement would have allowed us to assess, without interfering in any way with the second particle, that its momentum p2 = p-p1 is an element of reality. This would have allowed to conclude that p2 is an element of reality. Therefore, Einstein sums up, Quantum Mechanics is incomplete.
Bohr’s answer stresses once more that one cannot speak of quantities existing independently of the actual procedure of measuring them: "From our point of view we now see that the wording of the above mentioned criterion of physical reality proposed by EPR contains an ambiguity as regards the meaning of the expression “without in any way disturbing a system”. Of course there is, in a case like that just considered, no question of a mechanical disturbance of the system under investigation during the last critical stage of the measuring procedure. But even at this stage there is essentially the question of an influence on the very conditions which define the possible types of predictions regarding the future behaviour of the system. Since these conditions constitute an inherent element of the description of any phenomenon to which the term “physical reality” can be properly attached, we see that the argumentation of the mentioned authors does not justify their conclusion that quantum-mechanical description is essentially incomplete."
Einstein recognized that Bohr might be right, but remained attached to his own point of view (Bohr 1958b): “"To believe [that it should offer an exhaustive description of the individual phenomena] is logically possible without contradiction; - he admits - but it is so very contrary to my scientific instinct that I cannot forego the search for a more complete conception."
The question remained open for almost 50 years but was solved by two fundamental contributions. In 1964 John Bell (Bell 1964) showed that Einstein’s hypothesis of the existence of hidden variables capable of describing reality in more detail than QM might lead to an experimental test. In order to sketch Bell’s argument a reformulation of the original EPR proposal is necessary. Instead of chosing the relative distance and the total momentum as variables of the two-particle system with assigned initial value one assumes that they are two spin ½ particles in a state (singlet) of total angular momentum zero. In this state the components along three orthogonal directions are all zero, in spite of the fact that the three components of angular momentum are incompatible variable between themselves.
Bell’s idea is the following. Rather than discussing the legitimacy of speaking of a physical variable without having measured it, he proposes of measuring the component of the spin of particle #1 in a direction a and the component ofthe spin of particle #2 in another direction b. After a series of measurements on a great number N of pairs the results are correlated by a function C(a,b)=∑aibi (ai and bi may have values +/- 1) which depends on the angle θ between a et b. The point is, Bell shows, that Einstein’s hypothesis of hidden variables leads to an inequality
which is violated by the function C(a,b) = -cosθ of QM.
Bell’s inequality shows that the difference between Einstein’s and Bohr’s views is not only a matter of interpretation, but that the formalism of QM contradicts the hypothesis that incompatible variables may have at the same time sharp, even if unknown, values. The debate between Bohr and Einstein has been settled in favour of Bohr by Alain Aspect and coworkers (Aspect 1982) who showed in a celebrated experiment that the inequality (3) is violated for (ab)= 22,5 et (ab')= 67,5 by 5 standard déviations. Numerous other experiments have since then confirmed this result.
2. From quantons to objects
2.1. The existence of a classical world
We come back now to the ambiguous nature attributed by Bohr to the measuring apparatus. Does it belong to the classical or to the quantum world? In order to answer to this question we must preliminarly discuss the issue of the classical limit of Quantum Mechanics. We know that in the standard formulation of QM a system’s state is represented by a wave function in the cohordinate’s space (or a state vector in Hilbert space) which contains all the statistical properties of the system’s variables. The wave function allows to calculate the probability of finding a given value of any variable of the system as a result of a measurement by means of a suitable instrument. More precisely, if the wave function is given by
(where ψ1 (ψ2) represents a state in which the variable G has with certainty the value g1 (g2)), the probability of finding g1 (g2) is | c1|2 (|c2|2). In Bohr’s interpretation this means that the variable G does not have one of these values before its measurement but assumes one or the other value with the corresponding probability during the act of measurement. Now comes the question: is this interpretation always valid, even when g1 and g2 are macroscopically different?
The answer poses a serious problem. One can in effect prove that in the limit when Planck’s constant h tends to zero the probability distribution of the quantum state represented by ψ tends to the probability distribution in phase space of the corresponding classical statistical ensemble labeled by the same values of the system’s quantum variables. More precisely, in this limit |c1|2 and |c2|2 represent the probabilities of finding the values g1 (g2) of the classical variable corresponding to the quantum variable G. In this case, however, the interpretation of these probabilities is completely different. In classical statistical mechanics we assume that they express an incomplete knowlwdge of the values of G actually possessed by the different systems of the ensemble. We assume in fact that, if the ensemble is made of N systems, there are N|c1|2 systems wuth the value g1 of G and N|c2|2 systems with the value g2 of G to start with. Each system has a given value of G from the beginning, even if we don’t know it.
We arrive therefore to a contradiction. The same mathematical expression represents on the one side (classical limit of QM) the probability that a given system of the ensemble acquires a given value of the variable G as a consequence of its interaction with a suitable measuring instrument, and on the other side (classical statistical mechanics) the probability that the system considered had that value of G before its measurement. Suppose, for exemple that y represents a quantum in a box with two communicating cpmpartments: Ψ1 is different from zero in the left side compartment and Ψ2 is different from zero in the right side one. The corresponding probabilities of finding the quanton in one or the other one are respectively |c1|2 and |c2|2. Suppose now that the two compartments are separated by a shutter and displaced far away from each other. One of them is then opened: It may contain the quanton or may be empty. At this point it is undoubtably troubling to admit that, if one sticks strictly to Bohr’s interpretation, the system in question instantly materializes in one or the other locality when the compartment is opened. Even more troubling is the fact that, if QM is the only true and universally valid theory of matter, the same conclusion must hold in principle also for macroscopic bodies.
2.2. Quantum and classical uncertainties
A way out of this dilemma, however, exists. We have shown, with Maurizio Serva (Cini M. Serva M. 1990, 1992), that, without changing the basic principlees and the predictions of Quantum Mechanics, one can save at the same time both Bohr’s interpretation of the phenomena of the quantum domain, and Einstein’s belief in the objective relity of the classical world in which we live. We have shown in fact that the uncertainty product between x and p can be written for any state of a quanton in the form
where (∆x ∆p)q is of the order of the minimum value h/4π of the Heisenberg uncertainty relation and (∆x ∆p)cl is the classical expression of the product of the indeterminations ∆x and ∆p predicted by the probability distribution of the classical statistical mechanics distribution corresponding to the quantum state when h -> 0. It is therefore reasonable to attribute to each of these two terms the meaning relevant to its physical domain.
In the typical quantum domain the clssical term vanishes and the indeterminacy is ontological, namely the variables x and p do not have a definite value before the system’s interaction with a measuring instrument. When the accuracy of the act of measurement reduces the indeterminacy of one variable, the the indeterminacy of the other one increases. Their product cannot become smaller than h/4π.
As soon as the uncertainty product calculated from the state y acquires a classical term (which survives in the limit h -> 0) the total indeterminacy becomes epistemic, namely it represents an incomplete knowledge of the value that the measured variable really had before being measured. In this case it is possible to measure the variables x and p in such a way as to reduce at the same time both ∆x et ∆p without violating any quantum principle. These measurements reduce simply our ignorance. There is no instantaneous localization of the quanton in coordinate or momentum space as a consequence of the interaction between system and instrument, because position and momentum (within the intrinsic quantum uncertainty) were already localized.
This solution solves therefore the contradiction between the different interpretations of the total uncertainty product, and allows a reconciliation of the two alternative conceptions of physical reality proposed by Einstein and Bohr. It saves a realistic conception of the world as a whole by recognizing that macroscopic objects have objective properties independently of their being observed by any “observer”, and, at the same time, that at the microscopic objects have properties dependent of the macroscopic objects with which they interact,
It allows also to clarify the ambiguity on the nature of the measurment apparatus mentioned above. One can in fact reformulate it in the following way. Assume that the microscopic system S interacts with a part M1 which at its turn interacts with a part M2 and eventually other ones. We ask: at which point we pass the border between quantum domain and classical domain? The answer is not ambiguous. The border is where the values of the variable in one-to-one correspondence with the values of the quantum variable G, assume values which differ by each other by macroscopic quantities (e.g. charged or discharged counter). The part Mc when this happens is then the “pointer” of the instrument on whose unambiguous results all human observers agree.
This approch solves also a problem on which thousands of pages have been written, namely the problem of the “wave packet reduction” or “collapse” as a consequence of the act of measurement (Cini M, Levy Leblond J.M. 1991)(Wheeler J, Zurek W 1983). We recall that with this expression we mean that, after having measured G on a system S whose state is represented by ψ (eq.(1)) the wave function changes abruptly and instantaneously to ψ1 or ψ2 accordingly to the result g1 or g2 of the measurement. This change cannot be represented by a Schrödinger evolution, but must be postulated as a result of an instantaneous, irreversible and random evolution extraneous to QM. According to our findings (Cini M. et al 1979, Cini M. 1983) this additional and arbitrary mechanism is not necessary.
In fact, onsider the simplest case S+M, in which M is a counter which has two macroscopically different states (charged or discharged) represented by two state vectors Φ1 and Φ2. The wave funtion Ω of the total system may be written
where we have assumed that the value g1 (g2) of the variable G of S is correlated with the charged (discharged) counter. The preceding discussion shows that, due to the macroscopic difference between Φ1 and Φ2, the total systen’s state is, for all practical purposes, equivalent to a Gibbs classical ensemble made of N|c1|2 systems in which each counter is charged and S has the value g1 of G and N|c2|2 systems in which each counter is discharged charged and S has the value g2 of G. The wave packet reduction is therefore no longer needed as an additional postulate, and no additional misterious agent (even less the “observer’s consciousness”) is required to explain it. It simply turns out to be a well known consequence of classical statistical mechanics.
2.3. EPR and conservation laws
A similar "realistic" approach can be adopted to discuss the third counterintuitive quantum phenomenon, the famous EPR "paradox", whose solution, after the numerous experiments confirming the violation of Bell's inequalities, can only be expressed by saying that Einstein was wrong in concluding that quantum mechanics is an incomplete theory.
Usually people ask: how is it possible that when the first particle of a pair initially having zero total angular momentum acquires in interaction with its filter a sharp value of a given component of its angular momentum, the far away particle comes to "know" that its own angular momentum component should acquire the same and opposite value? I do not think that a realistic interpretation of this counterintuitive behaviour can be "explained" by minimizing the difference with its classical counterpart, because this difference has its roots, in my opinion, in the "ontological" (or irreducible) - not "epistemical" (or due to imperfect knowledge) - nature of the randomness of quantum events. If this is the case, one has in fact to accept that physical laws do not formulate detailed prescriptions, enforced by concrete physical entities, about what must happen in the world, but only provide constraints and express prohibitions about what may happen. Random events just happen, provided they comply to these constraints and do not violate these prohibitions.
From this point of view, the angular momentum component of the far away particle has to be equal and opposite to the measured value of the first particle's component, because otherwise the law of conservation of angular momentum would be violated. In fact, the quantity "total angular momentum" is itself, by definition, a non-local quantity. Non locality therefore needs not to be enforced by a mysterious action-at-a-distance. The two filters are not two uncorrelated pieces of matter: they are two rigidly connected parts of one single piece of matter which "measures" this quantity. The non local constraint is therefore provided by the nature of the macroscopic "instrument". This entails that, once the quantum randomness has produced the first partial sharp result, there is no freedom left for the result of the final stage of the interaction: there is no source of angular momentum available to produce any other result except the equal and opposite sharp value needed to add up to zero for the total momentum.
We arrive to the conclusion that Bohr was right, but Einstein was not wrong in insisting that an uncritical acceptance of the current interpretation of QM would lead to absurd statements about the physical nature of the world we live in.
3. The randomness of quantum reality in phase space
3.1. The representation of the irreducible randomness of quantum world in phase space
After eighty years of Quantum Mechanics (QM) we have learned to live with wave functions without worrying about their physical nature. This attitude is certainly justified by the extraordinary success of the theory in predicting and explaining not only all the phenomena encountered in the domain of microphysics, but also some spectacular nonclassical macroscopic behaviours of matter. Nevertheless one cannot ignore that the wave–particle duality of quantum objects not only still raises conceptual problems among the members of the small community of physicists who are still interested in the foundations of our basic theory of matter, but also induces thousands and thousands of physics students all around the world to ask each year, at their first impact with Quantum Mechanics, embarassing questions to their teachers without receiving really convincing answers.
We have seen that typical examples of this insatisfaction are the nonseparable character of long distance correlated two-particle systems and the dubious meaning of the superposition of state vectors of measuring instruments, and in general of all macroscopic objects (Schrödinger 1935). In the former case experiments have definitely established that Einstein was wrong in claiming that QM has to be completed by introducing extra “hidden” variables, but have shed no light on the nature of the entangled two-particle state vector responsible for the peculiar quantum correlation between them, a correlation which exceeds the classical one expected from the constraints of conservation laws.
In the latter case, generations of theoretical physicists in neoplatonist mood have insisted in claiming that the realistic aspect of macroscopic objects is only an illusion valid For All Practical Purposes (in jargon FAPP). The common core of their views is the belief that the only entity existing behind any object, be it small or large, is its wave function, which rules the random occurrence of the object’s potential physical properties. The most extravagant and bold version of this approach is undoubtedly the one known as the Many Worlds Interpretation of QM Everett E.(1973), which goes a step further by eliminating the very founding stone on which QM has been built, namely the essential randomness of quantum events. Chance disappears: the evolution of the whole Universe is written – a curious revival of Laplace - in the deterministic evolution of its wave function. “The Many-Worlds Interpretation (MWI) – in the words of Lev Vaidman, one of its most eminent supporters (Vaidman 2007) - is an approach to quantum mechanics according to which, in addition to the world we are aware of directly, there are many other similar worlds which exist in parallel at the same time and in the same space. The existence of the other worlds makes it possible to remove randomness and action at a distance from quantum theory and thus from all physics.”
I believe that it is grossly misleading to attribute the epistemological status of “consistent physical theory” to this sort of science fiction, which postulates the existence of myriads and myriads of physical objects (indeed entire worlds!) which are in principle undetectable. My purpose is to show that these difficulties can only be faced by pursuing a line of research which goes in the opposite direction, namely which takes for granted the irreducible nature of randomness in the quantum world. This can be done by eliminating from the beginning the unphysical concept of wave function. I believe that this elimination is conceptually similar to the elimination of the aether, together with its paradoxical properties, from classical electrodynamics, accomplished by relativity theory. In our case the lesson sounds: No wave funtions, no problems about their physical nature.
Furthermore, the adoption of a statistical approach from the beginning for the description of the physical properties of quantum systems sounds methodologically better founded than the conventional ad hoc hybrid procedure of starting with the determination of a system’s wave function of unspecified nature followed by a “hand made” construction of the probability distributions of its.physical variables. If randomness has an irreducible origin in the quantum world its fundamental laws should allow for the occurrence of different events under equal conditions. The language of probability, suitably adapted to take into account all the relevant constraints, seems therefore to be the only language capable of expressing this fundamental role of chance.
The proper framework in which a solution of the conceptual problems discussed above should be looked for is, after all, the birthplace of the quantum of action, namely phase space. It is of course clear that standard positive joint probabilities for both position and momentum having sharp given values cannot exist in phase space, because they would contradict the uncertainty principle. Wigner however, in order to represent Quantum Mechanics in phase space, introduced the functions called after his name (Wigner 1932) as pseudoprobabilities which may assume also negative values, and showed that by means of them one can compute any physically meaningful statistical property of quantum states.
A step further along this direction was made by Feynman (Feynman 1987), who has shown that, by dropping the assumption that the predictions of Quantum Mechanics can only be formulated by means of nonnegative probabilities, one can avoid the use of probability amplitudes, namely waves, in quantum mechanics. After all to the old questions about the physical meaning of probability amplitudes remains unanswered. Dirac said once “Nobody has ever seen quantum mechanical waves: only particles are detectable. Feynman is reported to have stated "It is safe to say that no one understands Quantum Mechanics". It is undeniable in fact that probability amplitudes are source of conceptual troubles (nonlocality of particle states, superposition of macroscopic objects' states).
The difficulty of introducing directly standard positive probability amplitudes in phase space in quantum mechanics arises, as is well known, from the impossibility of assigning precise values to incompatible variables. No joint probability density of x and p exists in phase space. However, negative probabilities - argues Feynman - have a physical interpretation.
"The idea of negative numbers - he writes - is an exceedingly fruitful mathematical invention. Today a person who balks at making a calculation in this way is considered backward or ignorant, or to have some kind of mental block. It is the purpose of this paper to point out that we have a similar strong block against negative probabilities. By discussing a number of examples, I hope to show that they are entirely rational of course, and that their use simplifies calculations and thought in a number of calculations in physics."
"If a physical theory for calculating probabilities yields a negative probability for a given situation under certain assumed conditions, we need not conclude the theory is incorrect. Two other possibilities of interpretation exist. One is that the conditions (for example, initial conditions) may not be capable of being realized in the physical world. The other possibility is that the situation for which the probability appears to be negative is one that can not be verified directly. A combination of these two, limitation of verifiability and freedom in initial conditions, may also be a solution to the apparent difficulty."
Admittedly, as he recognizes, a "strong mental block" against this extention of the probability concept is widespread. Once this has been overcome, however, the road is open for a new reformulation of Quantum Mechanics, in which the concept of probability “waves” is eliminated from the beginning. After all, particles and waves do not stand on the same footing as far as their practical detection is concerned. We have already remarked that the position of a particle assumes a sharp value as a consequence of a single interaction with a suitable detector, but we need a beam of particles to infer the sharp value of their common momentum. This means that we never detect waves: we only infer their existence by detecting a large number of particles.
A striking exemple of the usefulness of this approach is that the troubles of entangled states disappear. In fact the Wigner pseutoprobability of the singlet state of the EPR paradox is the product of the Wigner pseudoprobabilities of the two spin ½ particles. This means no more questions about the “superluminal transmission” of information between them.
3.2. Classical ensembles with “Uncertainty Principle”
Feynman's program, however, is still based on the conventional formalism of QM: state vectors in Hilbert space or wave functions in coordinates' space. In fact, Wigner's function W(q,p) (pseudoprobability density for sharp values q, p of incompatible variables q and p) is defined by the expression
which contains explicitly the wave function of the state..In Feynman’sapproach waves are therefore still needed to start with, because pseudoprobabilities are first expressed in terms of wave functions, and then forgotten. We will show, however, that it is possible to express Quantum Mechanics from first principles in terms of pseudoprobabilities without ever introducing the concept of probability amplitudes. This program has been recently carried on [Cini 1999] by generalizing the formalism of classical statistical mechanics in phase space with the introduction of two postulates (uncertainty and discreteness), which impose mathematical constraints on the set of quantum variables in terms of which any physical quantity can be expressed. QM is therefore reformulated in terms of expectation values of quantum variables as a generalization of the correspondent classical varibles of classical statistical mechanics, with the introduction of a single quantum postulate.
This goal will be attained in two steps. The first step is the formulation of a classical Uncertainty Principle. We consider all the classical ensembles of particles in phase space with coordinate q and momentum p in which a given variable A(q,p) has a well determined value α and its conjugate variable B(q,p) is completely undetermined -. Only ensembles of this kind in fact are the classical limit of the quantum states.
Following Moyal (1946), we will represent all the statistical properties of our ensembles, usually expressed by the joint probability distribution Pα(q,p), in terms of the expectation value Cα(k,x) (represented from now onwards by <....>α) of the "characteristic variable" C(k,x) = e(-i/h)(kq+xp) as follows
The requirement that all its systems have the value α of the variable A
entails that Cα(k,x) must satisfy the equation
where a(k,x) is the double Fourier transform of the function A(q,p).
Actually, eq. (8) is only apparently an integral equation, because it is easily reduced in terns of thr variables A and B to a simple algebraic functional equation with solution
In fact Pα(q,p) must be independent of B if this variable is indetermined in the ensemble. All this may seem trivial but actually it is not. Eq. (8) will be in fact one of our starting equations for the transition to QM.
We impose now that the result (9) should be invariant under the canonical transformations generated by any arbitrary function L
Therefore the Poisson Bracket of A wiyh L must satisfy
from which it follows that the characteristic function must satisfiy, in addition to (9), also the equation
for all k,x.
Eqs. (8) (12) are the formal expression of a “classical uncertainty principle”, representing the conditions to be fulfilled by classical ensembles having the property, invariant under canonical trasformations, that a given variable A has the value α and its conjugate variable B is undetermined. Up to now we are still in the domain of classical statistical mechanics.
3.3 The quantum postulate
The second, essential, step is to introduce the quantum into this scheme. This is done by imposing the fulfilment of a second postulate, based on the assumption that the founding stone of quantum theory is the experimental fact that physical quantities exist (the action of periodic motions, the angular momentum, the energy of bound systems..) whose possible values form a discrete set, invariant under canonical transformations, characteristic of each variable in question. This means that we should request that α belongs to a discrete spectrum independent of the phase space variables.
This feature can only be ensured if eq. (8) for the classical characteristic function Cα(k,x), which yields a continuous spectrum α for the values of the classical variable A, is modified to become a true Fredholm homogeneous integral equation for the quantum characteristic function C i(k,y) with a nonseparable kernel g(ky-hx), allowing for the existence of a discrete set of eigenvalues αi.
Similarly, eq.(12) expressing the uncertainty principle between the classical variables A and Β should be changed into
for the quantum characteristic function C i(k,x) of the ensemble caracterized by one of the values αi of the quantum variable A and by the complete indeterminacy of its quantum conjugate variable B. The functions g( ) and f( ) should be determined by imposing new self consistent rules for the quantum variables involved.
The two eqs (13) (14), however, cannot be obtained from (7) and (11) as in the classical case by ordinary commuting numbers.In fact the only way to obtain (13) (14) is to replace the classical characteristic variables C(k,x) obeying the standard rule of multiplication of exponentials with quantum variables C(k,x) having the property
and to replace their classical Poisson bracket with the Quantum Poisson Bracket
This means that, if we want to allow for the existence of discrete values of at least one variable L we are forced to represent all the variables A by means of noncommuting Dirac q-numbers.This means that the mathematical nature of the entities needed to represent the quantum variables is a consequence of the physical assumption of the discreteness of quantum variables and not viceversa, as the conventional view of reality underlying the conventional axiomatic formulation of Quantum Mecchanics assumes.
With (15) (16) the functions f( ) and g( ) turn out to have the expressions
As expected, the quantum variables C(k,x) with the properties (15) (16) turn out to have the same exponential form of classical statistical mechanics where the classical variables q and p are replaced by quantum variables q and p satisfying the commutation relations
of the standard variables of Quantum Mechanics
From the solution of equations (13) (14) one immediately obtains (by simple Fourier transform) the pseudoprobability Wi(q,p) corresponding to the quantum caracteristic function C i(q,p) of the ensemble. This pseudoprobability coincides with the Wigner function obtained from the standard wave QM wave function of the state. It is important to mention that all pseudoprobabilities satisfy the condition
which expresses the uncertainty principle in the reformulation of quantum theory in phase space. It is remarkable that this principle is given by an equality, thus eliminating the ambiguity of the Heisenberg inequality due to the presence of the two physically different terms appearing in eq. (1)
3.4 Field quantization in phase space and wave/particle duality
These results however leaved some conceptual problems still open. First of all, once the Schrödinger waves have been eliminated from Quantum Mechanics, how does one generalize its principles to Quantum Field Theory? One should not forget that, historically, QED was invented by Dirac (Dirac 1927) by submitting "first quantized" Schrödinger amplitudes to the procedure of "second quantization". If no "first quantized" probability amplitudes exist any more how does one proceed? And, secondly, isn't one throwing away the baby with the dirty water by forgetting that after all a quantum field must still show some of the wavelike properties of its classical limit?
A second paper [Cini 2003] has been therefore devoted to answer to these questions, leading to the conclusion that:
one should not start from nonrelativistic quantum mechanics in order to formulate quantum field theory, but viceversa;
the wavelike behaviour of the quanta of a quantum field is, as already Pascual Jordan had understood in 1926 [Born, Heisernberg, Jordan 1926], a straightforward consequence of imposing the Einstein property of discreteness to the intensity of a classical field - clearly a nonlocal physical entity - which exists objectively in ordinary three dimensional space.
It is appropriate to recall that for Jordan, in fact, it is quantization which brings into existence particles, both photons and electrons. According to him, therefore, rather than trying to explain phenomena like diffraction and interference of single particles as properties of "probability waves" one should simply view them as primary properties of the field of which they represent the quanta. "These considerations show - we read in his paper “On waves and corpuscles in quantum mechanics” [Jordan 1927] - that the quantized field is equivalent, in all its physical properties and especially with respect to its inensity fluctuations, to a corpuscular system (with a symmetric eigenfunction)".
The derivation of Wigner functions from the principles of uncertainty and discreteness illustrated in the previous paragraph provides the formalism for deducing the kind of wave/particle duality suggested by Jordan (and forgotten by the physicist's community since then) by simply imposing Einstein's quantization to the states of a classical field represented by means of statistical ensembles in the phase spaces of its normal modes.
Following the procedure sketched in the previous paragraph, we introduce a classical statistical ensemble for the r-th radiation oscillator of the field's normal modes defined by the constraint that the intensity Nr(q,p) has with certainty a given value νr. The equations (9) (11) remain valid, provided the variable A with its value α is replaced by the intensity N with its value ν and the conjugated variable B is replaced by the corresponding phase θ of each normal mode (we omit from now onwards the index r). Our procedure of field quantization will be based on the Einstein assumption of the existence of discrete field quanta. More precisely we assume that the spectrum of the quantum variable N of each field oscillator should be discrete. Eqs. (15) (16) remain unchanged and express now the result that, the quantum variables should be represented by means of non commuting quantities (Dirac's q-numbers). Quantization is therefore now a consequence of the physical property of the existence of field quanta, and not viceversa.
The field's states with a given number of quanta can now be represented by going from the quantum variables q, p to the Dirac complex variables a, a* expressed in terms of each wave's intensity N and phase θ by means of their standard expressions
The eigenvalue equations (13) (14) can be rewritten for the characteristic functions C n(β, β*) expressed in terms of the new variables β,β* related to k,x and h,y by means of the same relations (20). These equations can be solved to give the eigenvalues νn of the quantum variable N and their characteristic functions C n(β, β*) yielding
This result is expected, but remarkable, because it has been obtained by solving our new integral equations without any reference to Schrödinger wavefunctions. It is also easy with this formalism to treat the field's coherent states, as well as the processes of emission and absorption of photons from a source to reproduce the results obtained by Dirac in his seminal paper on the foundations of quantum electrodynamics. It turns out of course that the absorption rate is proportional to nr and the emission rate to nr+1 (Einstein's laws)
3.5. Conclusions
The main result of the reversal of the order of quantization from non relativistic quantum mechanics to quantum field theory gives a clear physical foundation to the mathematical nature of all quantum variables. The basic formal rules of quantum mechanics follow in this way from the Einstein postulate of the existence of field's quanta. The main conceptual result of this approach is therefore the clarification of the basic notion of wave/particle duality, which follows from this postulate, and simply reflects the dual nature of the quantum field as a unique physical entity objectively existing in ordinary three dimensional space (or ordinary four dimensional relativistic space, when is the case). From Jordan's point of view, in fact, the wavelike behaviour of any field's state with any number of discrete quanta simply reflects the property of a physical nonlocal entity which exists objectively in ordinary three dimensional space.
This goal has been achieved by imposing two requirements to the characteristic function (Moyal 1949) of the classical ensembles of thr field’s normal modes. The first one is that the probability distribution of the ensembles should be invariant under canonical transformations. The second requirement is quantization.
These two requirements are a reformulation of the principles introduced in the preceding nonrelativistic formulation of quantum mechnics where it was shown that the Wigner functions of the states of the one dimensional motion of a single particle can be directly derived without ever introducing Schrödinger wave functions. They lead to the two equations (18) and (21) whose solutions yield directly the quantum characteristic functions of the states of each mode, which turn out to be the double Fourier transforms of their Wigner functions. In the derivation of these equations one discovers that the field variables cannot be represented by ordinary numbers but should be represented by means of noncommuting mathematical objects.
With the direct construction of the Wigner functions of the states of quantum fields, the deBroglie-Schrödinger waves are thus eliminated from the formulation of quantum field theory. This means that, once that their nature of mathematical auxiliary tools has been recognized, the endless discussions about their queer physical properties, such as the nature of long distance EPR correlations between two or more particles or the meaning of the superposition of macroscopic states, become meanigless as those about the queer properties of the aether after its elimination declared by the theory of relativity.
Furthermore it supports the view that the most adequate representation of the random character of quantum phenomena ought to be based on Wigner-Feynman pseudoprobabilities in phase space, in which the constraints of the uncertainty principle are embodied, rather than insisting in representing them as events occurring in different spaces, (e.g. configuration or momentum) ruled by their correlated but separate classical probability laws. This view still meets a widespread resistance on the grounds that pseudoprobabilities are not positive definite, but is starting to acquire consensus in some domains of physics such as quantum optics (Leibfried et al 1988)leading even to a proposal for their experimental determination (Lutterbach et al1997).
Finally, the direct deduction of Wigner functions from first principles solves a puzzling unanswered question which has been worrying all the beginners approaching the study of our fundamental theory of matter, all along its 75 years of life, namely "Why should one take the modulus squared of a wave amplitude in order to obtain the corresponding probability?" We can now say that there is no longer need of an answer, because there is no longer need to ask the question.
Notes
- In what follows the variables are written in boldface and their values are in ordinary typeset. | https://www.intechopen.com/books/theoretical-concepts-of-quantum-mechanics/the-physical-nature-of-wave-particle-duality |
Atanu Dey has a plan to make all of India literate in 3 years because “for India, the most important infrastructure project is the one that will build its human capital base.”
First, the government of India must credibly commit to paying every literate and numerate person Rs 5,000 (about US$100). Second, ensure that every person who wants to learn basic literacy and numeracy can do so without having to pay a single penny. Third, provide testing centers around the country (especially in rural areas) where a person can be certified to have achieved basic literacy and numeracy. Finally, sit back and let the free market grind out the outcome which is total literacy within three years.
The details of this proposal follow from elementary logic and basic common sense. First, the cost-benefit analysis. There is long term cost of having about 300 million illiterate citizens. Each year, a literate person must be at least 10 percent more productive than an illiterate person. Assuming a per capita annual product of the illiterate population to be $200 (which is about half the annual per capita GDP of India), a 10 percent increase in productivity would be an increase of $20 per year per capita. Over a working life of about 40 years, that is an $800 increase in productivity per capita. Assume that the average working life of the 300 million illiterates of India is a conservative 20 years. Then the increase in additional product due to the additional 300 million literates is a conservative $120 billion (300 million times $20 times 20 years) in net present value terms.
I am using very conservative estimates of the benefits to make the case that the cost of doing so is a very small compared to the benefits. Assume very liberal costs of delivering basic literacy, say, $100 per capita. I will argue elsewhere that this is a very liberal estimate. Add to it $100, the incentive amount paid to the person upon passing a standardized test, and you have a total cost of $200 per capita. For the total population, it is amounts to $60 billion. This is half the aggregate social benefit estimated above.
Now one may ask, how will the government, which is totally inept as evidenced by the fact that 300 million Indians are illiterate despite lofty goals of making education univerally available and has not been able to make a dent even after over 57 years of spending huge amounts, be able to do this? The answer is simple: the government must not be in the business of providing the means and method of primary education. The only job of the government should be to finance the education. Let the private sector do the actual provisioning of education. | https://emergic.org/2004/10/01/educating-indias-300-million-illiterates/ |
Q:
If I wanted to generate a sequence of elements using the element position in the sequence as a variable, how would that variable be written?
I'm trying to create a sequence of elements which are dependent on their position in the sequence. What I mean is that I want the value $r$ in position 1 equal to 1, on position 2 equal to 2 and so on. E.g.
$$\langle r, r, r, r, r, r\rangle=\langle 1, 2, 3, 4, 5, 6\rangle$$
$$\langle 2^r, 2^r, 2^r, 2^r, 2^r, 2^r\rangle=\langle 2, 4, 8, 16, 32, 64\rangle$$
My first thought is having 2 sequences of elements. One of $r$ being all natural numbers and the second being the sequence I'm using $r$ in (The sequence called $\Bbb X$.) next to each other like so:
$$r=\langle 0, 1, 2, 3,\dots\rangle=\Bbb N$$
$$\Bbb X=\langle 2^{r+1}, 2^{r+1}, 2^{r+1}, 2^{r+1}, 2^{r+1}, 2^{r+1}\rangle$$
Both combining to make:
$$\langle 2, 4, 8, 16, 32, 64\rangle$$
Is there some sort of pre-existing symbol or notation to represent the value of whatever position an element is in a sequence, changing the value of an expression like $2^{r+1}$ based on its position?
A:
I will use curly brackets $\{\}$ to denote sets and angle brackets $\langle\rangle$ to denote sequences.
A set is something fundamentally different from a sequence, since sets are not ordered. That is, the sets $\{1,1,1,2,4,3\}$ and $\{4,2,3,1\}$ are equal: multiplicity and order of elements does not matter. For sequences the order does matter, so $\langle 1,1,2\rangle$, $\langle 1,2\rangle$ and $\langle 2,1\rangle$ are all distinct from each other.
If you want to create a sequence $\langle a_1,a_2,\dots,a_i,\dots\rangle$ (finite or infinite) with each element having a value dependent on its position, I would write it as $\langle f(i)\mid i\in I\rangle$, where $I$ is an ordered list of indices, and $f$ is a function with the index as input.
For example:
$\mathbb X=\langle 2^r\mid 1\leq r\leq 6\rangle=\langle2,4,8,16,32,64\rangle$
$\langle n\mid n\in\mathbb Z_{>0}\rangle=\langle1,2,3,4,\dots\rangle$
In both of these examples it is left implicit that the indices $r$ and $n$ are ordered from small to large in the usual way.
For creating sets you could use the same notation, but here the order doesn't matter. So as a set we have $\{2^r\mid 1\leq r\leq 6\}=\{2,4,8,16,32,64\}=\{64,2,8, 16,32,4\}$.
| |
The map shows where conference attendees came from. We’ve included the U.S. and Canada, but didn’t have room to show the home countries of attendees from Japan, Denmark, Netherlands, and Ireland.
391: Total attendance. Includes HPEN consumers (64 adults and 20 children); family members/caregivers (88 adult and 23 children); clinicians (45 dietitians; 11 nurses; 10 physicians; 8 pharmacists; 16 other); and exhibitors (106).
Videos and slide presentations of these talks and others are available to view, free of charge at: www.oley.org/2018confdoc. | https://oley.org/page/2018ConfOverview |
Each paper will be 5-6 pages, typed, double-spaced, in 12-point font with 1” margins, and should be formatted in proper MLA style.
Each research paper will include 4 recent outside sources. This means you will actually need to use the research in your paper, not simply list it on a works cited or bibliography page. Essays from the textbook do not count. Please use sources published within the last 5 years.
At least 3 of your sources must be from an academic source (a journal or scholarly book).
You must have proper in-text citations and a works cited page in MLA format. See the Clark Library website for help citing sources (Links to an external site.).
Your approach to your topic must be intersectional.
______________________________________________________________________________
Step #1: Choose a Topic
It is essential that you choose a topic that will hold your interest for the term, since your work will be focused around this topic for some time. Here are some suggestions to get you started. Please keep in mind that these are all broad categories and will need to be narrowed considerable to fit into a 5 page paper. So, use these suggestions as a starting point, and figure out from there what specific question or problem you’d like to research. You are also welcome to choose your own topic, as long as it fits with the parameters of the course. Here’s a tutorial that will help you brainstorm and focus your topic: http://libraryguides.library.clark.edu/brainstorming (Links to an external site.)
Choose a film or genre of films (2-3 films of the same type) and do a feminist analysis of those films.
reproductive justice
family leave polices
media representation of …
size discrimination
gendered violence
politics of housework
feminization of poverty
Environmental racism
internet activism
Step #2: Write the Proposal
Your proposal paper will be 1-2 pages discussing your topic, how you will focus your topic, what specific resources you intend to use, and why you chose this topic. You must do some preliminary research before writing this proposal to make sure you are able to find enough material (including academic sources) to meet the requirements of this assignment. Include at least one academic source that you have found during your preliminary research. For more details about how to find the sources you’ll need for this paper, see Step #3, below. DUE AT THE END OF WEEK 3
Step #3: Complete the Research
Don’t wait until the last minute to do so! The library has several different academic databases (EBSCO, ProQuest, Gale), which are a good place to look for academic sources. Make sure you check the “scholarly/peer reviewed” box when you are searching to find the appropriate materials. You should also check “full text” to return results that are full articles only. You can also access the library from home using your student ID/login. I highly recommend asking the librarian for help if you are not doing well in your search. You can also ask the librarian about inter-library loans, but make sure to do this early on so the material arrives in time for you to complete your project. Remember you will need a total of 4 resources, and at least three must be from scholarly sources. Also make sure that you record all of the source information, including author’s name, title of the article, title of the larger work, pages, editors, publisher, where it was published, website addresses, the date you looked at the material (online), etc. If you have a source without all of this information, you won’t be able to use it in your paper unless you are able to find it again, so write everything down as you are doing your research. Ask me or a librarian for help if you are confused about any of these instructions!
Step #4: Write the Paper
The final paper will be 5-6 pages. Make sure to follow the guidelines given above very carefully. Read these directions carefully and pay attention to all due dates. I strongly suggest students make use of campus resources, such as online tutoring (Links to an external site.). You can can submit your draft online to receive feedback! | https://qualityessays.net/research-paper-60/ |
The sea berry (sea buckthorn) bare root plants that were planted this year at the farm are looking really good. Only a couple of them didn’t survive. The farm has rather sandy soils and this was a great year for rain. The seaberry plants growing at the homestead look unhealthy and only half of the six survived. The soil here has more clay, is black in color and actually has worms/life in it.
Hearing how easy it is to propagate sea berry from cuttings I took cuttings from the best looking plants. It is late October so the cuttings will be indoors until the greenhouse is built.
I started researching how to get the sea berry seeds to germinate since none of the hundred plus I planted grew.
Mistakes I made include:
- planting in the spring (they need to be cold stratified (chiled) for 90 days)
- after they have been cold stratisfied they should be soaked in water for two days
- planting to deep (they should be exposed to sun light)
So this week I took some seeds out of the fridge (cold stratification). They have been in there for more than 90 days at this point. Soaked them in water for 2 days. Then placed them in zip lock bags with some wet sand from a stream (hearing that sand and water from a stream can help germination). Included a paper towel in one bag which helps to see the seeds and made sure the bags remain moist. Set these bags on a counter top exposed to the sun and on the second day some of the seeds have already germinated!! I moved some of the germinated seeds to the pots with the sea berry cuttings since that is the sandy soil they grew well in this year.
Update mid December.
Looks like most of the cuttings are alive. Some have green leaves on top and another looks like it will soon produce leaves many places. The seedlings are still growing and are now developing a tougher central stock. They are about two inches tall. | http://regenfarms.com/tag/propagating-sea-buckthorn/ |
Q:
Paper.js animating : Movement according to a paths normal
I am trying to animate a line and its normal ( another line ). But when i change position after or before setting rotation of the normal strange animation occurs.
Is there anybody who has an idea on that?
I have this code in sketchpad:
http://sketch.paperjs.org/#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
var outerH = 200;
var outerW = 300;
var group = new Group();
var spine = new Path({x:0, y:0});
spine.add({x:0, y:outerH/4});
spine.add({x:-outerW, y:outerH});
spine.strokeColor = 'red';
var nP = new Path();
nP.strokeColor = 'blue';
nP.add(new Point(0, 0))
nP.add(new Point(50, 0));
//nP.pivot = nP.bounds.topLeft;
group.addChildren([spine, nP]);
group.position = {x:200, y:300};
var loc = spine.getLocationAt(120);
var normal = spine.getNormalAt(120);
nP.position = loc.point;
nP.rotate(normal.angle);
view.onFrame = function(event) {
var sinus = Math.sin(event.time );
var cosinus = Math.cos(event.time );
// Change the x position of the segment point;
spine.segments[2].point.y += cosinus ;
spine.segments[2].point.x += sinus ;
}
If I uncomment -> nP.rotate(normal.angle); nP is not rotating with the line normal point?
A:
Please read the following post on the mailing list that explains this behavior and offers an option to switch paper.js into a mode that simplifies this scenario:
https://groups.google.com/forum/#!searchin/paperjs/applymatrix/paperjs/4EIRSGzcaUI/seKoNT-PSpwJ
| |
Guilin Seven Star Park is situated on the east bank of Li River, 1.5 kilometers away from the downtown of Guilin. Covering an area of 1.347 square kilometers, it is the largest comprehensive park in Guilin.
Seven Star Park got its name from the Seven Star Mountain inside it whose seven peaks bear a stunning resemblance to the Big Dipper in the sky. The park combines natural scenery with human landscapes, and it also includes a zoo. It gathers mountains, waters, caves, stones, pavilions, architecture and historical relics together, all of which make the Seven Star Park a masterpiece of parks in Guilin. The best visit time of this park is from April to October.
|What to See/Do:||Seven Star Mountain: The Seven Star Mountain has seven mountain peaks looking like the Big Dipper. This mountain is famous for its numerous rocks and peculiar caves including a Seven Star Cave. There is also a stone forest, a bonsai garden and over 500 stone inscriptions produced during Sui and Tang Dynasties.
|
Flower Bridge (Huaqiao): Originally built in the Song Dynasty, the Flower Bridge is the oldest bridge in Guilin. It is located near the frontgate of the Park and is totally 135 meters long. During spring and autumn, the two sides of the bridge are filled with blooming peach blossom and green bamboos, hence its name “Flower Bridge”.
Guihai Forest of Stone Tablets: There are over 220 pieces of stone inscriptions here; the inscriptions cover the contents of economy, military, culture, folk custom and so on. The stone inscriptions are presented to tourists in the forms of poem, essay, song, couplet, picture and more. The writings on the stones are various, such as regular script, cursive script, clerical script, and seal character etc.
|Duration:||2 hours|
|Admission:||75RMB (Several sites inside it require additional charges.)|
|Address:||No. 1 Qixing Road, Qixing District, Guilin, Guangxi Zhuang Autonomous Region.|
|Opening time:||From April to December: 06:00 – 19:30;
|
From November to March: 06:30 – 19:00. | http://www.guidewetravel.com/seven-star-park-qixing-gongyuan/ |
The end of the year is a hectic time. Kids are checking out as spring fever kicks in and you, as a teacher, are trying to keep it all together while providing your students with some different, extra motivating activities to help them focus.
This letter writing activity is one that fits the bill perfectly! The added bonus is not only is it a fun end of the year activity for your students, but it is a meaningful one that promotes gratitude and spreads kindness at your school.
I started doing this letter writing activity quite a few years ago, as I was wracking my brain, trying to find some end of the year activities that were a bit out of the norm. Once I had students writing these letters, I knew it was something that would become one of my end of the year teaching traditions.
So, you may be wondering what exactly is the letter writing activity, and how do you organize it? I’ll give you step by step instructions here!
1. Getting Started
To introduce the activity, I tell the kids that each of them will be writing a letter that he/she will deliver to an adult who works at our school. I make sure to tell the kids that the person they choose should be someone that they appreciate. It might be a custodian, yard duty, para, cafeteria worker, principal, office staff member, or even another teacher.
I do not let them write this letter to me, although each year kids ask me if that would be alright. For this project, I tell them that it needs to be someone from school but NOT me and that I would love it if they wrote a letter to me later on! The idea, besides reviewing friendly letter form, is to allow the child to let another person know how he/she is making a difference and that the child is grateful for what that person does.
2. Review Friendly Letter Format
I usually teach 4th or 5th graders, so friendly letter form is something most of them understand and just a quick review is all that is needed to jog their memories. If you teach younger students, you might need to spend a little more time teaching the friendly letter format. There are several different formats now for the friendly letter, (traditional, block with indents, and block without indents) so you can choose whichever one you are required to use or whichever one you like if you’re not bound to one in particular (I have all 3 formats included in this free set).
One thing I like to do when we review is to have all of the kids get out their whiteboards (we keep these in our desks) and do some “scribble writing” to show me what they think a friendly letter’s shape should look like. This 2 minute assessment gives me a good idea of how much we’ll need to review.
After that, I like to go step by step and show or make a ‘fake” but fun letter to someone (the principal, a popular singer/actor/athlete…) and we name the parts of the friendly letter as we go… the heading, the greeting, the body, the closing, the signature. Once we go over this, we do a scribble writing again on the whiteboards to make sure that everyone is on the right track for the letter.
One thing I do have the kids change from a typical letter is their address. I have everyone use the school’s address as their address for two reasons. One is that you might be surprised how many kids don’t know their exact address, and the second reason is just to protect the child’s address for privacy reasons, since I don’t send permission notes home to parents for the activity.
Besides the letter’s format, we do brainstorm how to start, and what kinds of things we might say in the letter. I always tell them that specific stories and remembrances are a great thing to add and I stress the idea that we need to take our time and to put some thought into the letter.
We also talk about the fact that this is a real letter, and not just an assignment to be turned in and we discuss how important the letters will be to the people reading them. This is a really good time to emphasize quality work.
3. Finishing the Letters
Besides writing the letters, I like kids to personalize them even more with an illustration. I show the kids how to make a simple border and add to it (flowers, geometric shapes, favorite things…). Instead of a border, they may also draw a picture on the bottom, the back or on a separate page. I also give them an envelope and show them where to write the person’s name and their name on it and how to fold the letter to make it fit nicely in the envelope. The kids may also decorate the envelope if they’d like to.
4. Delivering the Letters
Once each child’s letter is completed, I let each student deliver the letter to the recipient. I try to choose times to send them that aren’t disruptive to that person’s schedule or to our own.
A few minutes before recess or lunch is a good time to deliver the letters usually. If a recipient is not there on our delivery day, the child walks the letter to the office and the staff there puts these into the person’s box, to find when he/she returns.
I really love this activity for so many reasons. I love that it gently “forces”, I mean “encourages” kids to think about how grateful they are for the people in their school life and that they hopefully realize how expressing that gratitude means a lot to the person receiving it.
I also love how warm and fuzzy the letters make our school staff feel. Out of the blue, for a custodian, a librarian, or the lunch lady to find out that he/she is making a difference in the lives of kids is a beautiful thing!
If you’d like to download all of the forms I use for this KIT (Keep in Touch) letter writing activity, please visit my TpT store, where you’ll find them for free!
Looking for more end of the year activities? Here are two resources that are loaded with interactive fun! The End of the Year Activities (Set 1) actually includes this KIT activity and you’ll also love the End of the Year Activities (Set 2)
Also, if you need a memory book, editable awards, or even a 3rd – 5th grade Literacy Set for the End of the Year, I have one for 3rd – 5th grades that includes both print AND digital formats!
Click here to take a look at the End of the Year Print and Digital Memory Book.
Happy end of the year! | https://the-teacher-next-door.com/spread-some-kindness-letter-writing-activity/ |
Clinton Township Clintondale senior quarterback Dondray Paris sheds a tackle during a Sept. 14 game against St. Clair Shores South Lake.
Photo by Deb Jacques
CLINTON TOWNSHIP — In St. Clair Shores South Lake’s 40-0 Sept. 14 win over Clinton Township Clintondale, senior quarterback Josh Bogan had a first half so good that he barely broke a sweat in the last two quarters.
Bogan, a transfer from Detroit Martin Luther King, threw for 280 yards and four touchdowns — two to senior wide receiver Tyler Waters — all in the first half. Bogan finished with 310 yards passing. Waters finished with 183 receiving yards and an interception to go with the two scores.
“I came into this season trying to prove a point — that I could play,” Bogan said. “I just try to come out and help my team win. I don’t really care about the numbers. I just try to play as well as I can and help the team win.”
The Cavaliers with the win move to 3-1 overall, 2-1 Macomb Area Conference Silver Division. It’s one win more already than South Lake had all of last season. The 2017 season is the only season the Cavaliers have missed the playoffs under coach Vernard Snowden.
“We knew what we were getting into last season. We were young and inexperienced,” Snowden said. “We thought we’d be able to play better this season. To see it happen is good for the program.”
Clintondale falls to 1-3, 0-3 MAC Silver. The Dragons have dropped three straight for the first time since 2005.
“We just need more reps,” Clintondale coach Dave Schindler said. “We’ll come back and prepare and fight hard. We’ve got another game next week.”
The teams had some jitters early, trading punts on three straight possessions after a South Lake turnover on the game’s opening drive. The Cavs on their third drive found a groove, ending a quick three-play drive with a 27-yard touchdown pass from Bogan to Waters to take a 6-0 lead with 1:40 left the first.
Clintondale on its next drive came up empty after running 13 plays, including five with goal to go. South Lake on the first play of its next possession saw Bogan hit Waters for an 87-yard touchdown strike, as Waters made a grab just before the ball hit the turf before taking it all the way. A successful two-point try gave the Cavs a 14-0 advantage with 5:01 left in the half.
After a Clintondale three and out, South Lake would strike fast again. Bogan found running back Jerimyah Vines on a quick pass. Vines would do the rest, racing 76-yards up the sideline for a score to put the Cavs up 20-0 with 1:50 left in the half.
The South Lake defense continued its stellar play by forcing a fumble on Clintondale’s next possession. That set the stage for Bogan to continue his big night, as he found wide receiver Bruce Sammons on a fade route for an 11-yard touchdown to put the Cavs up 26-0 with 56.8 seconds left in the half.
“I don’t know what to say about (Bogan). He’s just a great kid, a great leader,” Snowden said. “He can make all of the throws.”
Similar to the start of the game, the second half began with punts on the first three possessions. South Lake continued its trend of quick strikes, as Vines took a handoff 50 yards for a touchdown to put the Cavs up 32-0 with 3:36 left in the third. South Lake went up 34-0 after another successful two-point try. Vines finished with 137 total yards and two touchdowns.
South Lake sophomore running back Jerrod Hennings got in on the action, scoring on a 15-yard run with 3:21 left to close out the scoring.
Snowden said Clintondale was able to keep South Lake’s run game in check, but the team was able to find a groove through the air.
“The guys in our secondary just couldn’t tackle,” Schindler said.
Clintondale junior running back Darell Walker led the Dragons with 100 rushing yards.
South Lake’s next game is against St. Clair Shores Lake Shore at 7 p.m. Sept. 21 at home.
Clintondale in Week 5 takes on Warren Woods Tower at 7 p.m. Sept. 21 on the road. | https://www.candgnews.com/news/qb-logan-has-a-big-night-as-south-lake-football-shuts-out-clintondale-109873 |
This application claims the benefit of U.S. Provisional Application No. 60/706,747 filed Aug. 10, 2005, the subject matter of which is hereby incorporated by reference.
1. Field of the Invention
The present invention relates generally to an application programming interface (API) for performing physics simulations. More particularly, the invention relates to an API for performing particle-based fluid simulations.
2. Description of the Related Art
An Application Programming Interface (API) is a set of definitions and/or protocols used to generate computer software applications (hereafter, applications). In general, an API defines how an application communicates with other software in a computer system. For example, an API may define how an application accesses (e.g., invokes, calls, references, modifies, etc.) a set of software components (e.g., functions, procedures, variables, data structures, etc.) in the system. Alternatively, an API may also define how an application interacts with a piece of software such as an interpreter (e.g., a JavaScript Interpreter).
An API typically takes the form of a set of “calls”, data structures, and variables that can be included in an application. The term “call” is used herein to denote any part of an application (e.g., an instruction, a code fragment, etc.) causing initiation, execution, retrieval, storage, indexing, update, etc. of another piece of software. In other words, including a particular call in an application generally causes the application to access a software component associated with the call.
The term “application” is used throughout this written description to describe any piece of software that enables at least one function in a computational platform. For example, an application typically comprises a data file that enables some function by providing a set of instructions for performing the function. The terms “routine” and “subroutine” are used to denote a part of an application smaller than the whole application. Each routine or subroutine in an application may comprise one or more subroutines and/or one or more calls to other software components as previously defined in this specification.
A classic example of a call included in an application is a “function call”. A function call generally comprises a name of a function and zero or more parameters or arguments for the function. When an application including the function call is executed, the function is invoked with its accompanying arguments.
The set of software components that can be accessed using a particular API is referred to as the implementation of the API. The implementation of an API may include, for example, a software library designed to run on a particular system. In general, an API is not bound to any one particular implementation. In fact, an application may be written using a particular API so that the application may be ported to systems with different implementations of the API. For instance, an API defined by the Open Graphics Library (OpenGL) standard allows programmers to write graphics applications that run on both UNIX and Windows based platforms, even though the underlying implementation of the API is different on each platform.
In general, the implementation of one API can be constructed with calls from another API. For example, an API defining complex high-level functions can be implemented using API calls defining intermediate-level functions. The API defining intermediate-level functions can be implemented using API calls defining low-level functions, and so on.
The implementation of most APIs is distributed in one of two ways. The first way is to include the implementation as part of a computer's operating system. For example, the implementation could comprise a set of code libraries distributed with the operating system. The second way is to distribute the implementation as a separate application or as an executable or a code library that has to be linked with and/or compiled with an application.
In some cases, source code for an API's implementation is available for viewing and modification. Where an API's implementation is available in this way, the API is called “open source”. In other cases, the API is only available as a set of binary files or the like, and hence, the only way to access the API is by including calls defined by the API in an application.
The term “software development kit” (SDK) is often used interchangeably with the term “API” because a SDK comprises a set of tools (e.g., an API) used to create software applications. However, a SDK generally comprises additional tools for developing applications besides an API. For instance, a software development kit may include utilities such as a debugger, or it may include special hardware tools used to communicate with embedded devices.
Most operating systems provide APIs that allow application programmers to create and/or control various system objects, such as display graphics, memory, file systems, processes, etc. In addition, operating systems may also provide APIs for performing common tasks such as multimedia processing, networking functions, and so forth.
Many independent software and/or hardware applications also provide APIs that allow application programmers to interface (e.g., control, communicate, etc.) with them. For example, an API used for communicating with a peripheral device such as a camera may define calls used to access low level software for adjusting specific aspects of the camera such as its aperture, shutter speed, exposure time, etc.
FIG. 1
FIG. 1
100
101
104
102
104
102
103
is a conceptual illustration of a software architecture for a system including an API. In , a plurality of applications access a plurality of software components (API routines) using an API . API routines form a part of the implementation of API and are located in an operating system or application .
FIG. 2
200
205
is a conceptual illustration of a conventional system used to run applications containing calls from an API . In this written description, the term “run” describes any process in which hardware resources associated with a computational platform perform an operation under the direction (directly or indirectly) of a software resource.
FIG. 2
FIG. 2
200
201
202
202
203
205
206
204
206
204
205
204
206
202
Referring to , system comprises a central processing unit (CPU) operatively connected to an external memory . External memory stores an application , API , and a plurality of API routines . Application invokes API routines using API calls defined by API . As indicated by broken boxes in , application and API routines run on CPU .
Using an API to write applications is advantageous for various reasons. For example, a good API generally provides a level of abstraction between an application programmer and the low level details of the software components called by the API. An API also generally provides access to commonly used functionalities (e.g., creating display graphics, formatting text, spawning processes) so that application programmers do not have to implement these functionalities from scratch. Moreover, since an API is not bound to a particular implementation, APIs generally allow applications to be ported between different systems. In addition, APIs provide standard representations for various programming tasks, which allows different programmers to work on an application without having to relearn all the calls contained therein.
Because of the various advantages provided by APIs, most specialized application areas in the field of computer science/engineering have associated APIs. The calls provided by an API serve as building blocks for creating applications of increasing complexity within the application area. For example, in computer graphics, low level APIs such as OpenGL and DirectX define functions for rendering primitive graphics components such as polygons and other shapes. Graphics programmers such as video game programmers then build upon the functions defined by the low level APIs to create higher level APIs defining functions for performing higher level tasks such as complex character animations. Accordingly, developing effective APIs increases the scope of what can be produced by application programmers within an application area.
An emerging application area where APIs are still in the primitive stages of development is the area of computational physics simulations. Computational physics simulations are used for a variety of purposes, ranging from scientific visualization to three dimensional (3D) game animation.
The goal of computational physics simulations is to model interactions between objects in a virtual world using the laws of physics. For example, in the case of scientific visualization, the physical forces and interactions between the elements of a polypeptide chain may be computationally modeled and observed in order to predict the conformation (e.g., folding) of a particular protein. In the case of 3D game animation, the physical forces and interactions between actors (e.g., characters, objects, substances, etc.) in a scene and their environment is modeled in order to generate lifelike animations of the scene. Simulations of forces such as gravity, pressure, friction, chemical forces, etc. can be combined to create lifelike animations of collisions, falling objects, explosions, and so forth.
FIG. 3
301
302
303
301
301
302
303
303
Formally defined, a “physics simulation” is a virtual representation of a physical entity or entities that changes over time in accordance with the laws of physics. For example, illustrates a physics simulation wherein a “world state” is periodically updated according to a step function to yield an “updated world state” . World state typically includes a set of objects having a number of associated physical attributes such as a size, shape, mass, location, velocity, acceleration, density, etc. World state also typically includes a set of forces acting on each of the objects. The forces may include, for example, gravity, pressure, friction, magnetic attraction/repulsion, etc. In step function , the objects are allowed to evolve, i.e., change physical attributes, for a predetermined time step in accordance with their associated velocities, forces, etc. A resulting new set of physical attributes and forces constitute updated world state . Step function can then be repeated to generate further world states from updated world state .
Examples of APIs designed to execute physics simulations include Havok Physics 3 (HP3). HP3 defines routines for performing collision detection, dynamics and constraint solving, and vehicle dynamics in a physics simulation.
One interesting problem that has not been adequately addressed by currently available APIs is the problem of simulating fluids (e.g., gases and liquids). A fluid simulation is a representation of a fluid that changes over time in accordance with the laws of physics. Because fluids may experience a host of complicated phenomena such as convection, diffusion, turbulence, and surface tension, simulating fluids can be difficult.
FIG. 4
A common approach to simulating fluids is to represent a fluid as a set of particles, each having an associated set of physical attributes such as size, mass, velocity, location, and a set of forces acting on it. illustrates a fluid simulation wherein the fluid is represented as a set of particles.
FIG. 4
401
402
403
Referring to , the set of physical attributes and forces associated with the particles constitute a fluid state . A step function acts on the particles for a predetermined time step so that the fluid will evolve in accordance with the velocities and forces associated with each particle. Accordingly, an updated fluid state is generated after the predetermined time step.
FIG. 4
403
405
404
405
In most practical scenarios, the fluid simulation is coupled to an output for displaying or evaluating properties of the simulated fluid. For example, in , updated fluid state is transformed into an image by a display function . Image could be generated, for example, by transforming the set of particles representing the fluid into a mesh representing a surface of the fluid and then rendering the mesh using standard graphics rendering hardware.
One particularly difficult problem relating to fluid simulations is being able to generate the fluid simulations in real-time. Real-time fluid simulations form an important part of many interactive computer environments such as medical simulations and video games. Unfortunately, conventional systems and methods have failed to produce satisfactory real-time fluid simulations because of the computational cost of simulating all of the forces on each particle and the lack of optimized hardware and software for performing the same.
Embodiment of the invention recognize the need for an API allowing programmers to incorporate realistic fluid simulations into their applications. The invention further recognizes the need for a system adapted to efficiently run applications written with the API.
In one embodiment, the invention provides a method of executing a physics simulation in a system comprising a computational platform, a main application stored in the computational platform, a secondary application stored in the computational platform, and a smoothed particle hydrodynamics (SPH) application programming interface (API) implemented in the computational platform. The method comprises defining a SPH call in the SPH API, by operation of the main application, invoking a software routine using the SPH call, and by operation of the secondary application, updating a state of the physics simulation in response to the software routine.
In another embodiment, the invention provides a system adapted to execute a physics simulation, the system comprising; a computational platform comprising at least one central processing unit and memory, a main application stored in the computational platform, a secondary application stored in the computational platform, a smoothed particle hydrodynamics (SPH) application programming interface (API) implemented in the computational platform, and a SPH call defined by the SPH API and included in the main application so as to invoke a software routine for updating a state of the physics simulation by operation of the secondary application.
The invention is described below in relation to several embodiments illustrated in the accompanying drawings. Throughout the drawings, like reference numbers indicate like exemplary elements, components, or steps. In the drawings:
FIG. 1
is a conceptual illustration of a software architecture for a system including an API;
FIG. 2
shows a conventional system adapted to run an application containing calls from a particular API;
FIG. 3
is an illustration of one example of a physics simulation;
FIG. 4
is an illustration of one example of a fluid simulation;
FIG. 5
shows a system adapted to run an application containing calls from a physics API in accordance with one embodiment of the present invention; and,
FIG. 6
is a flowchart illustrating a method of executing a physics simulation in accordance with one embodiment of the present invention.
The present invention recognizes the need for an API allowing programmers to incorporate realistic fluid simulations into their applications. The invention further recognizes the need for a system adapted to efficiently run applications written with the API.
The invention finds ready application in various application areas requiring lifelike fluid simulations, including, but not limited to application areas where the fluid simulations are computed in real-time. Exemplary application areas include, for example, computer games, medical simulations, scientific applications, and multimedia presentations such as computer animated films.
Exemplary embodiments of the invention are described below with reference to the corresponding drawings. These embodiments are presented as teaching examples. The actual scope of the invention is defined by the claims that follow.
FIG. 5
shows a system adapted to execute a physics simulation according to one embodiment of the invention.
FIG. 5
501
502
503
Referring to , the system comprises a computational platform including a first processor and a second processor operatively connected to each other and to an external memory . The term “computational platform” used herein refers broadly to any combination of computational hardware and memory used to process and store data. The term “processor” refers to a part of the computational platform (e.g., a logic circuit) used to process the data.
501
502
In one embodiment of the invention, first and second processors and comprise respective first and second CPUs operatively connected by a system bus. Accordingly, the first and second CPUs are generally located on a single printed circuit board (PCB).
501
502
In another embodiment of the invention, first and second processors and comprise a CPU and a co-processor operatively connected via a peripheral component interconnect (PCI) interface.
501
502
501
502
In yet another embodiment of the invention, first and second processors and comprise respective first and second cores of a dual core multiprocessor. In other words, first an second processors and are located on a single integrated circuit (IC) chip.
FIG. 5
Although shows two processors, the system could also comprise a single processor having at least two execution threads for executing multiple processes or applications in parallel.
FIG. 5
504
508
503
507
505
506
The system of further comprises a main application and a secondary application stored in an external memory . The system still further comprises an API defining a set of API calls used for invoking a corresponding set of API routines .
504
501
508
502
504
508
Main application typically runs substantially on first processor and secondary application typically runs substantially on second processor . Main application and secondary application typically run in parallel and use asynchronous signaling to communicate with each other.
504
508
FIG. 3
For example, in one embodiment of the invention, main application comprises a game program (i.e., a video game) and secondary application comprises a physics simulation of a virtual environment (i.e., a world state) of the game program. The game program typically spawns or initiates the physics simulation, including any necessary initialization. Thereafter, the physics simulation runs in parallel with the game program, updating the world state according to a predetermined time step as illustrated in . Whenever the world state is updated, the physics simulation sends a signal to the game program indicating a change in various physical attributes of the game program's virtual environment. Similarly, the game program also sends signals to the physics simulation whenever various actors in the game program interact with their environment in a way that might affect the physics simulation.
504
505
505
504
506
501
502
508
506
501
502
At least one of main application or secondary application contains one of API calls . API calls included in main application may cause API routines to run on first processor or on second processor . Likewise, API calls included in secondary application may cause API routines to run on first processor or second processor .
506
506
504
508
API routines are typically distributed as part of an operating system of the computational platform. However, API routines may also be distributed as an independent application or a set of independent executables or code libraries that have to be linked and/or compiled with main application and/or secondary application .
505
According to one embodiment of the invention, API calls are defined by a smoothed particle hydrodynamics (SPH) API. The SPH API defines calls for creating and executing particle based fluid simulations.
In general, the term “smoothed particle hydrodynamics” denotes a way of modeling a fluid using a finite set of particles. In smoothed particle hydrodynamics, various quantities related to the fluid (e.g., forces, mass, velocity, etc.) are defined at discrete locations in space (e.g., at the locations of corresponding particles). Then, the values the respective quantities are distributed over local neighborhoods of their various locations. In other words, the quantities are interpolated (i.e., “smoothed”) so that each quantity has a value defined at each location in space.
<math overflow="scroll"><mtable><mtr><mtd><mrow><mrow><mi>ρ</mi><mo></mo><mrow><mo>(</mo><mrow><mfrac><mrow><mo>∂</mo><mi>v</mi></mrow><mrow><mo>∂</mo><mi>t</mi></mrow></mfrac><mo>+</mo><mrow><mi>v</mi><mo>·</mo><mrow><mo>∇</mo><mi>v</mi></mrow></mrow></mrow><mo>)</mo></mrow></mrow><mo>=</mo><mrow><mrow><mo>-</mo><mrow><mo>∇</mo><mi>p</mi></mrow></mrow><mo>+</mo><mrow><mi>ρ</mi><mo></mo><mstyle><mtext> </mtext></mstyle><mo></mo><mi>g</mi></mrow><mo>+</mo><mrow><mi>μ</mi><mo></mo><mrow><mrow><msup><mo>∇</mo><mn>2</mn></msup><mo></mo><mi>v</mi></mrow><mo>.</mo></mrow></mrow></mrow></mrow></mtd><mtd><mrow><mo>(</mo><mn>1</mn><mo>)</mo></mrow></mtd></mtr></mtable></math>
<math overflow="scroll"><mtable><mtr><mtd><mrow><mrow><mrow><mi>ρ</mi><mo></mo><mrow><mo>(</mo><mfrac><mi>Dv</mi><mi>Dt</mi></mfrac><mo>)</mo></mrow></mrow><mo>=</mo><mrow><mrow><mo>-</mo><mrow><mo>∇</mo><mi>p</mi></mrow></mrow><mo>+</mo><mrow><mi>ρ</mi><mo></mo><mstyle><mtext> </mtext></mstyle><mo></mo><mi>g</mi></mrow><mo>+</mo><mrow><mi>μ</mi><mo></mo><mstyle><mtext> </mtext></mstyle><mo></mo><mrow><msup><mo>∇</mo><mn>2</mn></msup><mo></mo><mi>v</mi></mrow></mrow></mrow></mrow><mo>,</mo></mrow></mtd><mtd><mrow><mo>(</mo><mn>2</mn><mo>)</mo></mrow></mtd></mtr></mtable></math>
<math overflow="scroll"><mfrac><mi>Dv</mi><mi>Dt</mi></mfrac></math>
<math overflow="scroll"><mfrac><mi>Dv</mi><mi>Dt</mi></mfrac></math>
As an example of smoothed particle hydrodynamics, consider a set of particles defined at locations in a three-dimensional (3D) virtual world. Suppose that each particle has an associated mass, velocity, and acceleration. Further suppose that a plurality of forces (e.g., pressure and viscosity) and a density value are defined at the location of each of the particles. The values defined for each of the particles can be interpolated using a set of kernel functions such as radial basis functions centered at each of the particles' locations. Interpolating the values in this way generates continuous fields such as a pressure field and a viscosity field that can be used to compute the dynamics of the fluid using standard equations such as the Navier-Stokes equation. For example, one form of the Navier-Stokes equation models fluid dynamics as follows:
In equation (1), the term “v” denotes the velocity of a particle, the term “ρ” denotes a density value of the particle, the term “p” denotes a pressure force on the particle, the term “g” denotes gravity, and the term “μ” denotes a viscosity of the fluid. A convective term v·∇v in equation (1) can be ignored because it is assumed that each particle moves together with the fluid. Accordingly, the above form of the Navier-Stokes equation reduces to the following equation (2) for particle based physics simulations:
where the term
denotes the time derivative of the velocity of each particle. In other words,
denotes the acceleration of the particle. By interpolating the various quantities to generate continuous fields, the gradients of the velocity and pressure fields can be generated for equation (2) to allow the acceleration of each particle to be computed.
In addition to the pressure, viscosity, and gravity forces mentioned above, fluids can also be simulated using other forces such as surface tension. A method of simulating a fluid using smoothed particle hydrodynamics to simulate a fluid is described in further detail, for example, in “Particle-Based Fluid Simulation for Interactive Applications” (Matthias Müller et al., Eurographics/SIGGRAPH Symposium on Computer Animation, 2003).
Where kernel functions are used to interpolate the various quantities in smoothed particle hydrodynamics, the kernel functions may only be defined over a finite radius in order to limit the computational cost of interpolating the quantities. Alternatively, kernel functions with finite support may be used to interpolate the quantities. In either case, “a radius of influence” defines the range of the kernel function with non-zero values defined. For example, in one implementation of smoothed particle hydrodynamics, space is divided into a 3D grid of cells and the radius of influence for each particle inside a particular cell is defined to be the cell itself. In one exemplary embodiment of the SPH API, the radius of influence for particles is stored in a data structure “kernelRadiusMultiplier” that can be modified by an associated software routine.
Among the calls defined by the SPH API are calls for introducing particles into a particle based fluid simulation. Particles are typically introduced into the fluid simulation in one of two ways: by explicitly adding the particles, or by creating an emitter in the simulation.
Particles can be explicitly added into a simulation by specifying various properties of the particles such as locations, velocities, and positions, and then passing the properties to a software routine through a SPH API call so that the software routine will add the particles to the fluid simulation. For example, the SPH API call could comprise a function call “addParticles(const NxParticleData&)”, where “addParticles” is the name of a software routine for adding particles to the fluid simulation and “const NxParticleData&” is the type of a parameter denoting a reference to a data structure containing the locations, velocities, positions, and number of particles to be added to the fluid simulation.
Particles are added to a fluid simulation through an emitter by creating an emitter in the simulation that dispenses particles according to some predetermined pattern such as a rate. An emitter can be created by specifying various properties of the emitter such as size, location, pose, pressure, rate, etc., and passing the properties to a software routine through a SPH API call so that the software routine will create the emitter with the specified properties. For example, the SPH API call could comprise a function call “createEmitter(const NxFluidEmitterDesc&)”, where “createEmitter” is the name of a function for creating an emitter in the particle simulation and “const NxFluidEmitterDesc&” is the type of a parameter denoting a reference to a data structure containing the properties of the emitter.
The size of an emitter generally refers to the size of a cross sectional area across which particles enter the fluid simulation from the emitter. The pose of the emitter generally refers to the rotational orientation of the emitter, and the location of the emitter denotes the location of the emitter in 3D space.
As an alternative to, or in addition to, specifying properties of the emitter in the “createEmitter” call, various properties of the emitter can also be specified through other SPH API calls. For example, a function call “setGlobalPose(const NxMat34&)” in the SPH API could be used to invoke a software routine for setting the pose of an existing emitter in the fluid simulation. The function “setGlobalPose” receives a transformation matrix specified by parameter of a 3×4 matrix type “const NxMat34&” and applies the transformation to an emitter in order to rotate the emitter to a particular pose.
A function call “setGlobalPosition(const NxVec3&)” invokes a software routine for setting the location of an existing emitter to coordinates specified by parameter of a 3 dimensional vector type “const NxVec3&”.
A function call “setGlobalOrientation(const NxMat33&)” invokes a software routine for setting the orientation of an existing emitter according to a transformation matrix specified by parameter of a 3×3 matrix type “const NxMat33&”.
While the functions “setGlobalPose”, “setGlobalPosition”, and “setGlobalOrientation” establish the properties of the emitter in a “global” coordinate system, the SPH API also generally provides functions for modifying the pose, location, and orientation of the emitter in a local coordinate system. SPH API calls for modifying an emitter's properties in a local coordinate system may be used, for example, to place an emitter on a particular part of an actor within the fluid simulation. That way, properties such as the location, etc. of the emitter can remain fixed relative to the actor, but change relative to the world, e.g., by movement of the actor. Such functions could be invoked, for example by SPH API calls such as “setLocalPose(const NxMat34&)”, etc.
The rate of particles emitted by an emitter can also be established by a separate SPH API call. For example, a function call “setRate(NxReal)” invokes a software routine that sets a number of particles emitted by the emitter per unit of time, e.g, particles per second, according to a parameter of type “NxReal”.
Similarly, the velocity of particles produced by the emitter can also be set by a separate SPH API call. For example, a function call “setFluidVelocityMagnitude(NxReal)” invokes a software routine for setting a velocity of particles emitted by the emitter. The velocity is specified by a parameter of type “NxReal”.
An emitter can be removed from a simulation to stop the production of new particles by using yet another SPH API call. An example of a SPH API call for removing an emitter is a function call “releaseEmitter(NxFluidEmitter&)”. The function name “releaseEmitter” denotes a software routine for deleting an emitter from a fluid simulation and the parameter type “NxFluidEmitter&” denotes a reference to the emitter to be deleted.
Another property of particles in a physics simulation that can be set by a SPH API call is a lifetime of the particles. The particle's lifetime is an interval of time between when the particle was created, either by an explicit SPH API call or by an emitter, and when the particle is automatically deleted from the simulation. The lifetime of a particle can be set, for example, by a function call “setParticleLifetime(NxReal)”, which invokes a routine setting the lifetime of a particle to a value specified by a parameter of type “NxReal”.
In addition to the properties assigned to individual particles, the SPH API can also assign specific properties to all of the particles in the fluid simulation. For example, in one embodiment the SPH API contains calls for setting a stiffness, viscosity, number of particles per linear meter in a rest state, and a density of the fluid in the rest state.
The stiffness of particles is specified by a factor related to pressure. The factor linearly scales the force acting on particles which are closer to each other than a spacing of the particles in the rest state. A good value for the stiffness depends on many factors such as viscosity, damping, and the radius of influence of each particle. Values which are too high will result in an unstable simulation, whereas too low values will make the fluid appear “springy” (the fluid acts more compressible). The SPH API includes a function call “setStiffness(NxReal)” invoking a routine to set the stiffness factor for the simulation according to a parameter of type “NxReal”.
The viscosity of the fluid defines its viscous behavior. Higher values result in a honey-like behavior. Viscosity is an effect which depends on the relative velocity of neighboring particles; it reduces the magnitude of the relative velocity. The SPH API includes a function call “setViscosity(NxReal)” invoking a routine to set the viscosity for particles in the fluid simulation according to a parameter of type “NxReal”.
The number of particles per linear meter of the fluid in the rest state is the number of particles in a unit volume of the fluid when the net forces on each of the particles is zero. In other words, it is the number of particles that exist in a linear meter of the fluid when the particles are standing still and not acting on each other, i.e., in the “rest state”. The SPH API includes a data structure “restParticlesPerMeter” for specifying this parameter. The data structure can be updated by a software routine invoked by the SPH API or it can be directly modified by an application.
rd
The rest density of a fluid is the mass of the fluid per unit volume when the fluid is in the rest state. For example, the rest density of water is about 1000 kilogram per meter cubed at 4 degrees Celsius. This parameter defines indirectly, in combination with the number of particles per linear meter in the rest state, the mass of one particle, i.e. the mass of one particle=“rest density”/“particles per cubed meter in the rest state”. Particle mass has an impact on repulsion effect on emitters and actors. The repulsion effect is the force generated on the emitters and actors due to Newton's 3law of physics. The SPH API includes a data structure “restDensity” for specifying the rest density. This data structure can be updated by a software routine invoked by the SPH API or it can be directly modified by the application.
As described above, a fluid simulation may be executed by a main application and a secondary application running in parallel. In particular, the fluid simulation may be executed substantially by the secondary application and another program controlling the physics simulation may be executed by the main application. Accordingly, the SPH API provides a function call for receiving data from the fluid simulation. For example, the main application may include the function call for receiving the data from the secondary application in order to gather and manipulate data related to the fluid simulation. The SPH API includes a function call “setParticlesWriteData(const NxParticleData&)” invoking a routine for receiving data from the fluid simulation. The data is stored in a data structure indicated by a parameter of type “const NxParticleData&”.
The system described above is one example of a system that could execute a physics simulation using a SPH API. Other systems adapted to execute physics simulations using an SPH API are disclosed, for example, in U.S. patent applications with Ser. Nos. 10/715,459 and 10/715,440 filed Nov. 19, 2003, Ser. No. 10/815,721 filed Apr. 2, 2004, Ser. No. 10/839,155 filed May 6, 2004, Ser. No. 10/982,791 filed Nov. 8, 2004, and Ser. No. 10/988,588 filed Nov. 16, 2004.
FIG. 6
FIG. 6
is a flowchart illustrating a method of executing a physics simulation in a system such as those described above. In the description of , parentheses (XXX) are used to indicate exemplary method steps.
FIG. 6
601
602
603
Referring to , the physics simulation is executed by defining an SPH call in the SPH API (). Then, by operation of the main application, a software routine is invoked using the SPH call (). For example, the SPH call may be included in the main application so that when the main application executes, the SPH call causes a corresponding software routine to execute. Then, by operation of the secondary application, the physics simulation is updated in response to the software routine (). For example, the software routine could add or delete an emitter or update various properties of particles in a fluid simulation using any of various SPH API calls described above.
BACKGROUND OF THE INVENTION
SUMMARY OF THE INVENTION
BRIEF DESCRIPTION OF THE DRAWINGS
DESCRIPTION OF THE EXEMPLARY EMBODIMENTS | |
PROBLEM TO BE SOLVED: To provide a RUBI (KANA (Japanese syllabary) characters written alongside Chinese characters to indicate their pronunciation) processing method capable of arranging easily readable RUBI of excellent appearance and usability.
SOLUTION: A RUBI processing part 20 includes a document reading part 200 and a character acquiring part 201 for acquiring characters from a read document. A character total width calculation part 202 and a RUBI total width calculation part 203 respectively calculate the total width Cw of continuous characters having RUBI and the total width Rw of RUBI written alongside the continuous characters. An width comparing part 204 compares the Cw with the Rw, and when both the Cw and Rw coincide with each other, a closing processing part 205 closely arranges these characters and the RUBI without generating gaps. Thereby the characters and the RUBI are arranged in the same width.
COPYRIGHT: (C)1999,JPO | |
As a young girl, Sarah painted the finishing details in the decorative areas of her father’s canvases, like lacework and flowers. Her first public works date from 1816 with subjects such as flowers and still lifes, but she soon turned to portrait painting. At the age of eighteen, Sarah drew a self-portrait, a tradition in her family. If the portrait was successfully done, her family would consider her an artist instead of a student.
James Peale was not pleased enough with Sarah’s self-portrait, so she decided to do something different from the rest of her family – something other than still-lifes and miniatures. In 1818, she spent three months in Baltimore, Maryland studying with her cousin, noted portrait artist Rembrandt Peale, learning new techniques, and again in 1820 and 1822. He greatly influenced her painting style and subject matter.
After experimenting with still lifes and miniatures, Sarah Peale exhibited her first full-size portrait at the Pennsylvania Academy of the Fine Arts in 1818. In 1824, Sarah and her sister, miniaturist Anna Claypoole Peale, became members of the Pennsylvania Academy of Fine Arts, America’s most prestigious institute. They were the first women to achieve this distinction.
Sally [Sarah’s nickname] also possesses great talents, her first and second attempts in Portrait are now exhibiting in the Academy of the Fine Arts and each of the girls [Sarah and her sister Anna Claypoole Peale] have had their share of praise by the critics in the newspaper.
By 1820 Sarah was often occupied with portrait commissions in Baltimore. Starting in 1822 she exhibited there annually at the Peale Museum, and by 1825, she maintained a studio there. Her portraits are distinctive for their detailed furs, laces, and fabrics – and realistic skin, faces, and hair.
Sarah Peale was not only a great artist, she was a pioneer for single, independent professional women. She left her family home to live and work on her own. Although there was virtually no precedent for women pursuing art professionally in the United States, Sarah launched an independent career as a portrait artist, working in Philadelphia and in Baltimore, Maryland.
Peale opened her studio in Baltimore and established herself as one of Baltimore’s most capable portrait painters. At this time, prior to photography, there was a wide market for portraits, and she attracted some of the best subjects. Diplomats, congressmen, and other public figures wanted to be drawn by her, and then bought her paintings.
Some of the men she painted include Massachusetts Senator Daniel Webster, Missouri Senator Thomas Hart Benton, and French General Marquis de Lafayette. Sarah painted over one hundred portraits of members of Baltimore society, more than any of her competitors in that city.
For 25 years, she painted in Baltimore (1822–1847), with occasional trips to Washington, DC, where she attended sessions of Congress and painted portraits of men in the government. More than 100 commissioned portrait paintings are known from her time in Baltimore and she was the most prolific artist in the city during that era. Her subjects were wealthy Baltimore residents and politicians from Washington DC.
Miss Sarah M. Peale intends visiting our city the approaching fall, for the purpose of painting several portraits. Four specimens from the pencil of this lady are now in the Missouri bank, and they clearly prove her title to rank among the first of American artists.
Peale became independently successful in St. Louis and continued to earn a living through her work. She became successful as a painter of portraits, mainly of politicians and military figures, and the occasional still life, while working in St. Louis (1847-1878). There, she continued to be a leading portrait painter of her day.
Records show that Peale received many more portrait commissions than celebrated male painters of that time, such as Thomas Sully and John Vanderlyn. Most of her work from this era is in private hands and not available for viewing.
Around 1860 Peale inexplicably returned to painting still lifes, but with more natural arrangements than her earlier works. She won numerous awards for these works, as well.
Sarah Peale returned to Philadelphia in 1878, having lived away from home for more than fifty years, first in Baltimore, and then in St. Louis. She spent the last years of her life in Philadelphia, living with her sisters, Anna Claypoole Peale (1791-1878) and Margaretta Angelica Peale (1795-1882). Like her sisters, Sarah never married, preferring to devote her energies to her career.
Sarah Miriam Peale died February 4, 1885 at age 85, and was buried at the Gloria Dei (Old Swedes’) Church.
Sarah Miriam Peale is noted as a portrait painter, mainly of politicians and military figures. Marquis de Lafayette sat for her four times. She won numerous awards throughout her life, and she was one of the first women in the United States to achieve professional recognition as an artist. She maintained a career for more than fifty years and supported herself without marrying, which was almost unheard of in the mid-nineteenth century. | http://www.womenhistoryblog.com/2015/04/sarah-miriam-peale.html |
CROSS-REFERENCE TO RELATED APPLICATIONS
BACKGROUND OF THE INVENTION
SUMMARY OF THE INVENTION
BRIEF DESCRIPTION OF THE DRAWINGS
DETAILED DESCRIPTION OF THE INVENTION
This application claims the benefit of priority under 35USC §119 to Japanese Patent Application No. 2003-310368 filed on Sep. 2, 2003, No. 2004-19552 filed on Jan. 28, 2004, and No. 2004-233503 filed on Aug. 10, 2004, the entire contents of which are incorporated herein by reference.
1. Field of the Invention
The present invention relates to an inverse model calculation apparatus and an inverse model calculation method.
2. Background Art
It is one of problems demanded in the field of control or the like to find an input required to obtain a desirable output from an object system (inverse calculation). If physical characteristics of the object system are already obtained as a numerical expression, the input can be found by solving the numerical expression.
In many cases, however, the numerical expression is not obtained beforehand. In the case where the numerical expression is not obtained beforehand, typically a mathematical model representing characteristics of the object system is constructed by using data obtained by observing the object system.
Typically, a forward model used to find an output obtained when a certain input is given can be constructed easily. However, it is difficult to generate an inverse model used to find an input required to obtain a certain output. The reason is that there are a plurality of inputs for which the same output is obtained.
Therefore, it is frequently performed to first construct a forward model, and estimate an input from an output by using the forward model. In such a case, a method using a generalized inverse matrix of a linear model, a method of performing an inverse calculation using a neural net, a solution by using simulation, and so on have heretofore been used.
However, the method using the generalized inverse matrix of a linear model becomes poor in calculation precision in the case where the nonlinearity of the object system is strong or in the case of multi-input and a single output.
On the other hand, in the inverse calculation using a neural net, all input variables used to construct the forward model of the neural net become the calculation object, and consequently even an unnecessary input is identified, and it is difficult to find an optimum input. Furthermore, in the inverse calculation using the neural net, it is difficult to calculate after how many time units the given output is obtained.
The solution using simulation is a method of giving various inputs to a forward model and determining whether a target output is obtained in a cut and try manner. Therefore, a large quantity of calculation is needed, and consequently it takes a long time to perform the calculation.
In order to solve the above-described problem, the present invention provides an inverse model calculation apparatus and an inverse model calculation method capable of efficiently calculating an input condition required to obtain a desired output.
An inverse model calculation apparatus according to an embodiment of the present invention provides an inverse model calculation apparatus for finding a condition under which a target system outputs a certain output value, the target system outputting the certain output value on the basis of an input value to the target system, the inverse model calculation apparatus comprising: a time series data recording section which records an input value inputted sequentially to the target system and an output value outputted sequentially from the target system as time series data; a decision tree generation section which generates a decision tree for inferring an output value at future time, using the time series data; and a condition acquisition section which detects a leaf node having an output value at future time as a value of an object variable from the decision tree, and acquires a condition of explaining variables included in a rule associated with a path from a root node of the decision tree to the detected leaf node, as a condition for obtaining the output value.
An inverse model calculation apparatus according to an embodiment of the present invention provides an inverse model calculation apparatus for finding a condition under which a target system outputs a certain output value, the target system outputting the certain output value on the basis of an input value to the target system, the inverse model calculation apparatus comprising: a time series data recording section which records an input value inputted sequentially to the target system and an output value outputted sequentially from the target system as time series data; a decision tree generation section which generates a decision tree for inferring an output value at future time, using the time series data; a condition acquisition section which an output value at future time is inputted into as a initial condition, which detects a leaf node having the inputted output value as a value of an object variable from the decision tree, and which acquires a condition of explaining variables included in a rule associated with a path from a root node of the decision tree to the detected leaf node, as a condition to obtain the output value; and a condition decision section, which determines whether the acquired condition is a past condition or a future condition, which determines whether the acquired condition is true or false by using the time series data and the acquired condition in the case where the acquired condition is the past condition, which determines whether the acquired condition is an input condition or an output condition in the case where the acquired condition is the future condition, which outputs the acquired condition as a necessary condition for obtaining the output value in the case where the acquired condition is the input condition, and which outputs the acquired condition to the condition acquisition section as an output value at future time in the case where the acquired condition is the output condition.
An inverse model calculation apparatus according to an embodiment of the present invention provides an inverse model calculation apparatus for finding a condition under which a target system outputs a certain output value, the target system outputting the certain output value on the basis of an input value to the target system, the inverse model calculation apparatus comprising: time series data recording section which records an input value inputted sequentially to the target system and an output value outputted sequentially from the target system as time series data; a decision tree generation section which generates a decision tree for inferring an output value at future time, using the time series data, a path from a root node to a leaf node being associated in the decision tree with a rule including a condition of explaining variables and a value of an object variable; a first rule detection section which detects a rule having an output value at future time as a value of an object variable, from the decision tree; a first condition calculation section which determines whether a condition of explaining variables for a partial time zone in the detected rule matches the time series data, and which in the case of matching, calculates a condition for obtaining the output value at the future time, using the detected rule and the time series data;
a second rule detection section, to which a rule is inputted, and which detects a rule that a condition of explaining variables for a partial time zone in the inputted rule matches from the decision tree; a first input section which inputs the rule detected by the first rule detection section to the second rule detection section, in the case where the rule detected by the first rule detection section does not match the time series data; a second input section which determines whether a condition of explaining variables for a partial time zone in the rule detected by the second rule detection section matches the time series data, and which, in the case of not-matching, inputs the rule detected by the second rule detection section to the second rule detection section; and a second condition calculation section which calculates a condition for obtaining the output value at the future time, using all rules detected by the first and second rule detection sections and the time series data, in the case where the rule detected by the second rule detection section matches the time series data.
An inverse model calculation method according to an embodiment of the present invention provides an inverse model calculation method for finding a condition under which a target system outputs a certain output value, the target system outputting the certain output value on the basis of an input value to the target system, the inverse model calculation method comprising: recording an input value inputted sequentially to the target system and an output value outputted sequentially from the target system as time series data; generating a decision tree for inferring an output value at future time, using the time series data; and detecting a leaf node having an output value at future time as a value of an object variable from the decision tree; and acquiring a condition of explaining variables included in a rule associated with a path from a root node of the decision tree to the detected leaf node, as a condition for obtaining the output value.
An inverse model calculation method for finding a condition under which a target system outputs a certain output value, the target system outputting the certain output value on the basis of an input value to the target system, the inverse model calculation apparatus comprising: recording an input value inputted sequentially to the target system and an output value outputted sequentially from the target system as time series data; generating a decision tree for inferring an output value at future time, using the time series data; inputting an output value at future time as a initial condition; detecting a leaf node having the inputted output value as a value of an object variable from the decision tree; acquiring a condition of explaining variables included in a rule associated with a path from a root node of the decision tree to the detected leaf node, as a condition for obtaining the output value; determining whether the acquired condition is a past condition or a future condition; determining whether the acquired condition is true or false by using the time series data and the acquired condition in the case where the acquired condition is the past condition; determining whether the acquired condition is an input condition or an output condition in the case where the acquired condition is the future condition; outputting the acquired condition as a necessary condition for obtaining the output value in the case where the acquired condition is the input condition regarding the acquired condition as an output value at future time in the case where the acquired condition is an output condition, and detecting a leaf node having the regarded output value at the future time as a value of an object variable from the decision tree, acquiring a condition of explaining variables included in a rule associated with a path from the root node to the detected leaf node, as a condition for obtaining the regarded output value.
An inverse model calculation method for finding a condition under which a target system outputs a certain output value, the target system outputting the certain output value on the basis of an input value to the target system, the inverse model calculation method comprising: recording an input value inputted sequentially to the target system and an output value outputted sequentially from the target system as time series data; generating a decision tree for inferring an output value at future time, using the time series data, a path from a root node to a leaf node being associated in the decision tree with a rule including a condition of explaining variables and a value of an object variable; detecting a rule having an output value at future time as a value of an object variable, from the decision tree; in the case where a condition of explaining variables for a partial time zone in the detected rule matches the time series data, calculating a condition for obtaining the output value at the future time, using the detected rule and the time series data; in the case of non-matching, newly detecting a rule matching the condition of explaining variables for a partial time zone in the detected rule, from the decision tree; in the case where a condition of explaining variables for a partial time zone in the newly detected-rule does not match the time series data, further detecting a rule which the condition of explaining variables for a partial time zone in the newly detected rule matches, from the decision tree; repeating detecting a rule which a condition of explaining variables for a partial time zone in a latest detected rule matches, from the decision tree, until a rule whose condition of explaining variables for a partial time zone matches the time series data is detected; and calculating a condition required to obtain the output value at the future time by using all rules detected from the decision tree and the time series data, in the case where the rule whose condition of explaining variables for a partial time zone matches the time series data has been detected.
FIG. 1
is a block diagram showing a configuration of an inverse model calculation apparatus according to a first embodiment of the present invention.
FIG. 2
shows an input sequence of a variable X input to a target system and an output sequence of a variable Y output from the target system.
FIG. 3
1
2
is a diagram showing time series data including input sequence of variables X and X input to the target system and an output sequence of a variable Y output from the target system, in a table form.
FIG. 4
FIG. 3
is a diagram showing a decision tree generated on the basis of the time series data shown in .
FIG. 5
1
2
is a diagram showing time series data including input sequence of variables X and X and an output sequence of a variable Y in a table form.
FIG. 6
FIG. 5
1
2
is a table showing a data obtained by regarding the variable Y as an object variable and the variables X and X as explaining variables and rearranging the time series data shown in .
FIG. 7
is a flow chart showing processing steps performed by the inverse model calculation apparatus.
FIG. 8
is a flow chart showing processing steps of the subroutine A.
FIG. 9
is a block diagram showing a configuration of an inverse model calculation apparatus according to the second embodiment.
FIG. 10
FIG. 9
is a flow chart showing the processing steps performed by the inverse model calculation apparatus shown in .
FIG. 11
is a flow chart showing processing steps in the subroutine B.
FIG. 12
is a flow chart showing processing steps performed by the inverse model calculation apparatus according to the third embodiment of the present invention.
FIG. 13
FIG. 3
is a table showing a part that follows the time series data shown in .
FIG. 14
is a diagram showing time series data to be analyzed.
FIG. 15
FIG. 14
is a table showing a state in which the time series data shown in have been rearranged.
FIG. 16
FIG. 15
shows a decision tree constructed on the basis of the table shown in .
FIG. 17
is a diagram showing the rules (1) to (13) in a table form.
FIG. 18
is a diagram explaining the logical inference.
FIG. 19
is a diagram showing concretely how logical inference is performed by combining the rule (10) with the rule (4).
FIG. 20
is a flow chart showing processing steps performed by the inverse model calculation apparatus.
FIG. 21
is a flow chart showing processing steps in the subroutine C in detail.
FIG. 22
is a flow chart showing processing steps in the subroutine D.
FIG. 23
is a flow chart showing processing steps in the subroutine E.
FIG. 24
is a block diagram showing a configuration of an inverse model computer system using an inverse model calculation apparatus.
FIG. 25
is a configuration diagram of a decision tree combination apparatus, which combines a plurality of decision trees.
FIG. 26
shows another example of the decision tree combination apparatus,
FIG. 27
is a table showing an example of observed data.
FIG. 28
1
shows data used to generate one decision tree (a decision tree associated with the object variable Y).
FIG. 29
1
2
is a diagram showing examples of the decision tree and the decision tree .
FIG. 30
is a flow chart showing a processing procedure for performing the combination method 1.
FIG. 31
shows a example of a series of explaining variable values.
FIG. 32
shows one generated instance data.
FIG. 33
is a flow chart showing a processing procedure for performing a combination method 2.
FIG. 34
1011
is a flow chart showing a processing procedure at the step S.
FIG. 35
is a diagram showing an example of a path set.
FIG. 36
FIG. 35
is a diagram showing a state in which the path set shown in have been concatenated.
FIG. 37
FIG. 36
shows a path (composite path) obtained by eliminating the duplication from the concatenated path shown in .
FIG. 38
16
shows generated composite paths.
FIG. 39
1012
is a flow chart showing the processing procedure at the step S in detail.
FIG. 40
shows the decision tree in the middle of generation.
FIG. 41
shows the decision tree in the middle of generation.
FIG. 42
shows the decision tree in the middle of generation.
FIG. 43
shows the decision tree in the middle of generation.
44
1
2
FIG. shows a decision tree obtained by combining the decision tree with the decision tree .
FIG. 45
is a flow chart showing a processing procedure for performing a combination method 3.
46
FIG. shows the decision tree in the middle of generation.
47
FIG. shows the decision tree in the middle of generation.
FIG. 48
1
2
shows a decision tree obtained by combining the decision tree with the decision tree .
FIG. 49
is a diagram showing an evaluation method of a leftmost path in the composite decision tree.
(First Embodiment)
FIG. 1
8
is a block diagram showing a configuration of an inverse model calculation apparatus according to a first embodiment of the present invention.
1
1
1
A time series data recording section records input values inputted sequentially to a target system as an input sequence. A time series data recording section records output values outputted sequentially from the target system as an output sequence. A time series data recording section records the input sequence and the output sequence as time series data (observed data).
FIG. 2
4
4
shows an input sequence of a variable X input to a target system and an output sequence of a variable Y output from the target system .
FIG. 3
FIG. 3
1
2
4
4
4
is a diagram showing time series data including input sequence of variables X and X input to the target system and an output sequence of a variable Y output from the target system , in a table form. As shown in the , in this target system , a one-dimensional output sequence is output on the basis of a two-dimensional input sequence.
2
1
FIG. 1
A decision tree generation section shown in generates a decision tree for inferring an output sequence on the basis of an input sequence by using time series data stored in the time series data recording section .
FIG. 4
FIG. 3
is a diagram showing a decision tree generated on the basis of the time series data shown in .
1
1
2
1
2
4
1
In this decision tree, an output Y(t) at time t can be predicted on the basis of an input sequence of a variable X supplied until time t. Among the input sequence of the two variables X and X, only the input sequence of the variable X appears in this decision tree, and the input sequence of the variable X does not appear. In other words, in this target system , the output Y can be predicted from only the input sequence of the variable X. In this way, there is an effect of reducing the input variable used for the prediction by using a decision tree. The decision tree has a plurality of rules. Each rule corresponds to a path from a root node of the decision tree to a leaf node. In other words, the decision tree includes as many rules as the leaf nodes.
Here, as the specific generation method of the decision tree, an already known method can be used. Hereafter, the method required to generate the decision tree will be described briefly.
FIG. 5
1
2
is a diagram showing time series data including input sequence of variables X and X and an output sequence of a variable Y in a table form.
First, the already known method is applied to this time series data to rearrange this time series data.
FIG. 6
FIG. 5
1
2
is a table showing a data obtained by regarding the variable Y as an object variable and the variables X and X as explaining variables and rearranging the time series data shown in .
FIG. 6
Subsequently, a method described in “C4.5: Programs for Machine Learning,” written by J. Ross Quinlan, and published by Morgan Kaufmann Publishers, Inc., 1993 is applied to the data shown in . As a result, a decision tree for predicting the output on the basis of the input sequence can be generated.
FIG. 1
FIG. 4
3
2
10
3
3
1
10
1
8
3
3
3
Returning back to , a condition acquisition section acquires a condition required to obtain an output value at a given future time by tracing branches of the decision tree generated by the decision tree generation section from a leaf node toward the root node. For example, if an output Y()=3 is given as an output at a future time in , then the condition acquisition section specifies a leaf node corresponding to the output in the decision tree, traces branches from the leaf node to the root node, and detects X()>=2 and X()<1. In other words, the condition acquisition section specifies a rule having the output as the leaf node, and acquires a condition included in this rule as a condition required to obtain the output .
8
FIG. 1
Processing steps performed by the inverse model calculation apparatus shown in will now be described.
FIG. 7
8
is a flow chart showing processing steps performed by the inverse model calculation apparatus .
2
1
1
First, the decision tree generation section generates a decision tree by means of time series data recorded by the time series data recording section (step S).
3
2
Subsequently, an output value (Y(t)=V) (output condition) at a future time is given to the condition acquisition section by using data input means or the like, which is not illustrated (step S).
3
3
The condition acquisition section executes a subroutine A by regarding the output condition as a target condition (step S).
FIG. 8
is a flow chart showing processing steps of the subroutine A.
3
First, the condition acquisition section retrieves a leaf node having a target value (=V) in the decision tree (step
12
3
13
If there is no leaf nodes having the target value (NO at step S), then the condition acquisition section outputs a signal indicating that the condition required to obtain the target value cannot be retrieved, i.e., the target value cannot be obtained (false) (step S).
12
3
14
On the other hand, if there is a leaf node having the target value (YES at the step S), then the condition acquisition section traces the tree from the retrieved leaf node toward the root node, specifies a condition required to obtain the target value, and outputs the condition (step S).
100
FIG. 4
As a concrete example, it is now assumed that a condition required to obtain the target value 3 at time is to be retrieved by using the decision tree shown in .
FIG. 4
1
12
1
98
1
100
1
14
In the decision tree shown in , retrieval of a leaf node having a target value=3 is performed. As a result, a leaf node having the target value=3 is retrieved (Y(t)=3) (the step S and YES at the step S). Assuming that t=100, a condition that X()<1 and X()>=2 is obtained by tracing the tree from the leaf node to the root node (X(t)) (step S).
8
FIG. 1
An example of an inverse model computer system using the inverse model calculation apparatus shown in will now be described below.
FIG. 24
8
is a block diagram showing a configuration of an inverse model computer system using an inverse model calculation apparatus .
6
4
4
8
4
8
6
6
An input sequence generation section generates an input sequence of a variable X to be given to a target system . The target system generates an output sequence of a variable Y on the basis of the input sequence of the variable X. An inverse model calculation apparatus acquires the input sequence and the output sequence from the target system . The inverse model calculation apparatus implements the above-described processing, calculates an input condition required to obtain an output value at a given future time, and outputs the calculated input condition to the input sequence generation section . The input sequence generation section generates an input sequence in accordance with the input condition input thereto.
8
FIG. 1
FIG. 24
Heretofore, the inverse model calculation system incorporating the inverse model calculation apparatus shown in has been described. In the same way as the present embodiment, inverse model calculation apparatuses of the second to fifth embodiments described hereafter can also be incorporated in the inverse model computer system shown in .
According to the present embodiment, a decision tree is constructed as a model, and an input condition required to obtain an output value at a given future time is calculated, as heretofore described. Therefore, the amount of calculation can be reduced, and calculation of a value of an input variable that does not exert influence upon the output can be excluded.
According to the present embodiment, a decision tree is constructed as a model. Even if nonlinearity of the target system is strong, therefore, the precision of the model can be remained high.
(Second Embodiment)
The first embodiment shows a typical example of the inverse calculation using a decision tree, and it is indistinct whether the obtained condition can be actually satisfied. In the present embodiment, inverse calculation including a decision whether the obtained condition can be actually satisfied will now be described.
FIG. 9
is a block diagram showing a configuration of an inverse model calculation apparatus according to the second embodiment.
1
2
3
Since the time series data recording section , the decision tree generation section , and the condition acquisition section are the same as those of the first embodiment, detailed description thereof will be omitted.
3
5
3
5
If an output condition is included in conditions obtained by the condition acquisition section , then a condition decision section performs retrieval again by using the condition acquisition section and using the output condition as the target condition. The condition decision section repeats this processing until all conditions required to obtain a given output value are acquired as the input condition.
FIG. 9
Hereafter, processing steps performed by the inverse model calculation apparatus shown in will be described in detail.
FIG. 10
FIG. 9
is a flow chart showing the processing steps performed by the inverse model calculation apparatus shown in .
2
1
21
First, the decision tree generation section generates a decision tree by using time series data recorded by the time series data recording section (step S).
2
5
22
Subsequently, the decision tree generation section gives an output value at a future time (a target condition) to the condition decision section by using data input means, which is not illustrated (step S).
5
23
100
101
102
3
100
1
101
2
102
5
23
Subsequently, the condition decision section generates a target list, which stores the target condition (step S). The target list has a form such as “Y()=3, Y()=1, Y()=2, . . . ” (output at time , output at time and output at time ). On the other hand, the condition decision section especially prepares an input list, which stores obtained input conditions, and empties the input list (step S).
5
24
In this state, the condition decision section executes a subroutine B (step S).
FIG. 11
is a flow chart showing processing steps in the subroutine B.
5
31
First, the condition decision section determines whether the target list is empty (step S).
31
5
32
5
100
100
101
102
101
102
If the target list is not empty (NO at the step S), then the condition decision section takes out one of items from the target list (step S). For example, the condition decision section takes out the target condition “Y()=3” from the above-described target list “Y()=3, Y()=1, Y()=2, . . . .” In this case, items in the target list is decreased by one, resulting in “Y()=1, Y()=2, . . . .”
5
33
10
1
The condition decision section determines whether the item taken out is a past condition (step S). If the current time is provisionally , then a target condition “Y()=2” is a past condition.
33
5
34
5
If the item taken out is a past condition (YES at the step ), then the condition decision section determines by using past time series data whether the item taken out is true or false (step S). In other words, the condition decision section determines whether the item taken out satisfies past time series data.
34
5
35
If the decision result is false, i.e., the item taken out does not satisfy past time series data (false at the step S), then the condition decision section outputs a signal (false) indicating that the given output value cannot be obtained (step S).
34
5
31
On the other hand, if the decision result is true, i.e., the item taken out satisfies past time series data (true at the step S), then the condition decision section returns to the step S.
33
33
5
36
If it is found at the step S that the item taken out is not a past condition, i.e., the item taken out is a future condition (NO at the step S), then the condition decision section determines whether the item is an input condition or an output condition (step S).
36
5
3
37
5
3
100
5
3
100
5
3
FIG. 8
If the item taken out is an output condition (output condition at the step S), then the condition decision section causes the condition acquisition section to execute the subroutine A shown in by using the output condition as the target condition (step S). In other words, the condition decision section requests the condition acquisition section to retrieve a condition required to achieve that target condition. For example, if the item “Y()=3” taken out from the above-described target list is a future condition, then the condition decision section causes the condition acquisition section to execute the subroutine A by using “Y()=3” as the target condition. The condition decision section receives a retrieval result from the condition acquisition section .
3
38
5
35
If the retrieval result received from the condition acquisition section is false (YES at the step S), i.e., if a leaf node having the target value under the target condition is not present in the decision tree, then the condition decision section outputs a signal indicating that an output value at a given future time cannot be obtained (false) (step S).
3
38
3
5
39
On the other hand, if the retrieval result received from the condition acquisition section is not false (NO at the step S), i.e., if a condition (an input condition, an output condition, or an input condition and an output condition) required to achieve the target condition is received from the condition acquisition section as the retrieval result, then the condition decision section adds this condition to the target list as a target condition (step S).
36
36
5
40
1
100
1
101
2
100
If the item taken out is an input condition at the step S (input condition at the step S), then the condition decision section adds this input condition to the input list (step S). The input list has a form such as “X()=2, X()=3, X()=1 . . . .”
5
31
31
5
41
Thereafter, the condition decision section returns to the step S, and repeats the processing heretofore described. If the target list has become empty (YES at the step S), then the condition decision section outputs an input condition stored in the input list, as a necessary condition required to obtain an output value at a given future time (outputs true) (step S).
The present embodiment has been described heretofore. If the detected condition is a past condition, then truth or false is determined by comparing this condition with past time series data. If the detected condition is a future output condition, then retrieval is performed repeatedly. Therefore, it can be determined whether an output value at a given future time will be possible, if possible, a condition required to obtain the output value can be acquired as an input condition.
(Third Embodiment)
In the present embodiment, how many time units after the current time an output value at a given future time can be acquired at the shortest will now be described.
FIG. 9
5
A configuration of an inverse model calculation apparatus in the present embodiment is basically the same as that in the second embodiment shown in . However, the present embodiment differs from the second embodiment in processing performed by the condition decision section .
Hereafter, the inverse model calculation apparatus in the present embodiment will be described.
FIG. 12
is a flow chart showing processing steps performed by the inverse model calculation apparatus according to the third embodiment of the present invention.
2
1
51
First, the decision tree generation section generates a decision tree by using time series data recorded by the time series data recording section (step S).
2
5
52
Subsequently, the decision tree generation section gives an output value at a future time (supplies a target condition) to the condition decision section by using data input means, which is not illustrated (step S).
5
53
1
8
9
Subsequently, the condition decision section substitutes an initial value 0 for time t (step S). As for the initial value, the last time when an output value is present in the above-described time series data is substituted. (For example, if there are input values and output values at time to time and only an input value at time in the time series data, then the last time becomes 8.) Here, 0 is substituted as the initial value for brevity of description.
5
5
54
Subsequently, the condition decision section substitutes t+1 for time t. In other words, the condition decision section increases the time t by one (step S). This “1” is, for example, an input spacing time of the input sequence inputted to the target system.
5
55
Subsequently, the condition decision section determines whether the time t is greater than a predetermined value (step S).
55
5
56
If the time t is greater than the predetermined value (YES at the step S), then the condition decision section outputs a signal indicating that the given output value V cannot be obtained within the predetermined time (step S).
55
5
57
On the other hand, if the time t is equal to the predetermined value or less (NO at the step S), then the condition decision section empties the target list and the input list (step S), and adds a target condition “Y(t)=V” (output V at time t) to the target list.
5
59
FIG. 11
Upon adding the target condition “Y(t)=V” to the target list, the condition decision section executes the above-described subroutine B (see ) (step S).
60
5
54
55
59
If a result of the execution of the subroutine B is false (YES at step S), i.e., an input condition required to achieve Y(t)=V cannot be obtained, then the condition decision section further increases the time t by one (step S) and repeats the above-described processing (steps S to S).
60
5
61
On the other hand, if the result of the execution of the subroutine B is not false (NO at the step S), i.e., an input condition required to achieve Y(t)=V can be obtained, then the condition decision section outputs the input condition and the value of the time t (step S).
Processing steps performed by the inverse model calculation apparatus heretofore described will be further described by using a concrete example.
FIG. 13
FIG. 3
2
is a table showing a part that follows the time series data shown in . However, the variable X is omitted.
1
16
1
17
An input value of the variable X and an output value of the variable Y until time , and an input value of the variable X at time are already obtained.
An example in which the inverse model calculation apparatus calculates what time the output value subsequently becomes 3 (Y(t)=3) will now be described.
2
51
5
52
FIGS. 3 and 13
FIG. 4
FIG. 12
First, the decision tree generation section generates a decision tree by using time series data shown in (the same tree as that shown in is generated) (the step in ). Subsequently, a target condition (Y(t)=3) is input to the condition decision section via input means, which is not illustrated (step S).
5
16
53
5
The condition decision section substitutes for time t (step S). In other words, the condition decision section substitutes the last time when an output value exists for t.
5
17
54
The condition decision section increases the time t by one to obtain (step S).
5
55
5
55
57
The condition decision section determines whether the time (=17) is greater than a predetermined value (step S). Here, the condition decision section determines t (=17) to be equal to the predetermined value or less (NO at the step S), and empties the target list and the input list (step S).
5
17
58
59
5
60
The condition decision section adds a target condition “Y()=3” to the target list (step S), and executes the subroutine B (step S). The condition decision section determines the execution result to be false (YES at step S).
17
1
15
1
17
31
32
33
36
37
38
39
1
15
1
15
31
32
33
34
39
17
5
35
34
FIG. 4
FIG. 13
In other words, to achieve Y()=3 when t=18, it is necessary to satisfy X()<1 and X()>=2 as represented by the decision tree shown in (the steps S, S, S, S, S, NO at S, and S in the subroutine B). As shown in , however, X is 2 at time , and consequently the above-described X()<1 is not satisfied (steps S, S, S and false at step S following the step S). At the time , therefore, the condition decision section determines that the output value Y=3 cannot be obtained (step S following the step S).
5
54
18
5
57
58
59
5
60
FIG. 12
As a result, the condition decision section returns to the step S as shown in , and increases t by one to obtain . And the condition decision section executes the subroutine B again, via the steps S and (step S). Here as well, the condition decision section determines the execution result to be false (YES at step S).
18
1
16
1
18
31
32
33
36
37
38
39
1
16
1
16
31
32
33
34
39
18
5
35
34
FIG. 4
FIG. 13
In other words, to achieve Y()=3 when t=18, it is necessary to satisfy X()<1 and X()>=2 as represented by the decision tree shown in (the steps S, S, S, S, S, NO at S, and S in the subroutine B). As shown in , however, X is 3 at the time , and consequently X()<1 is not satisfied (steps S, S, S and false at step S following the step S). At the time , therefore, the condition decision section determines that the output value Y=3 cannot be obtained (step S following the step S).
5
54
19
5
57
58
59
5
60
FIG. 12
As a result, the condition decision section returns to the step S as shown in , and increases t by one to obtain . And the condition decision section executes the subroutine B again, via the steps S and (step S). Here as well, the condition decision section determines that the execution result is false (YES at step S).
19
1
17
1
19
31
32
33
36
37
38
39
1
17
1
17
31
32
33
34
39
19
5
35
34
FIG. 4
FIG. 11
FIG. 13
In other words, to achieve Y()=3 when t=19, it is necessary to satisfy X()<1 and X()>=2 as represented by the decision tree shown in (the steps S, S, S, S, S, NO at S, and S in ). As shown in , however, X is 3 at the time , and consequently X()<1 is not satisfied (steps S, S, S and false at step S following the step S). At the time , therefore, the condition decision section determines that the output value Y=3 cannot be obtained (step S following the step S).
5
54
20
5
57
58
59
5
60
FIG. 12
As a result, the condition decision section returns to the step S as shown in , and increases t by one to obtain . And the condition decision section executes the subroutine B again, via the steps S and (step S). The condition decision section determines that the execution result is not false (NO at step S).
20
1
18
1
20
31
32
33
536
37
38
539
531
532
33
539
5
36
540
533
5
20
31
41
540
60
61
FIG. 4
FIG. 11
FIG.12
In other words, to achieve Y()=3 when t=20, it is necessary to satisfy X()<1 and X()>=2 as represented by the decision tree shown in (the steps S, S, S, , S, NO at S, and in ). Both of these two input conditions are future conditions (steps , , NO at step S following ). Therefore, the condition decision section adds these two input conditions to the input list (input condition at step S and following the step ). The condition decision section outputs an input condition in the input list and the value of time t (YES at step S, S following the step , and NO at the step S and S in ).
According to the present embodiment, an input condition required to obtain a given output value is retrieved while successively increasing the value of a future time t as heretofore described. Therefore, it is possible to calculate how many time units after the current time an output value at a given future time can be acquired at the shortest.
(Fourth Embodiment)
In the present embodiment, an input condition required to obtain an output value at a given future time is calculated by performing “logical inference” using a plurality of rules (paths from the root node to leaf nodes) included in the decision tree and using time series data.
3
5
The present embodiment differs from the second and third embodiments in processing contents performed by the condition acquisition section and the condition decision section .
Hereafter, the present embodiment will be described in detail.
FIG. 14
is a diagram showing time series data to be analyzed.
2
1
2
The time series data are rearranged by regarding Y at time t as an object variable and regarding X at the time (t-) to t and Y at time t-, t- as explaining variables.
FIG. 15
FIG. 14
is a table showing a state in which the time series data shown in have been rearranged.
FIG. 16
FIG. 15
2
A decision tree is constructed by applying an already known method to this table. shows a decision tree constructed on the basis of the table shown in . This decision tree is generated by the decision tree generation section .
3
13
The condition acquisition section traces branches of this decision tree from the root node to a leaf node, and acquires the following rules (paths).
1
2
1
(1) Y(T−)<=4, Y(T−)<=5, X(T)=0, X(T−)=0→Y(T)=6
1
2
1
(2) Y(T−)<=4, Y(T−)<=5, X(T)=0, X(T−)=1→Y(T)=5
1
2
1
(3) Y(T−)<=4, Y(T−)<=5, X(T)=1, X(T−)=0→Y(T)=4
1
2
1
(4) Y(T−)<=4, Y(T−)<=5, X(T)=1, X(T−)=1→Y(T)=6
1
2
(5) Y(T−)<=4, Y(T−)=>6, X(T)=0,→Y(T)=5
1
2
1
(6) Y(T−)<=4, Y(T−)=>6, X(T)=1, X(T−)=0→Y(T)=5
1
2
1
(7) Y(T−)<=4, Y(T−)=>6, X(T)=1, X(T−)=1→Y(T)=6
1
2
2
(8) Y(T−)=>5, Y(T−)<=5, X(T)=0, X(T−)=0→Y(T)=4
1
2
2
(9) Y(T−)=>5, Y(T−)<=5, X(T)=0, X(T−)=1→Y(T)=5
1
2
(10) Y(T−)=>5, Y(T−)<=5, X(T)=1,→Y(T)=4
1
2
1
(11) Y(T−)=>5, Y(T−)=>6, X(T)=0, X(T−)=0→Y(T)=6
1
2
1
(12) Y(T−)=>5, Y(T−)=>6, X(T)=0, X(T−)=1→Y(T)=4
1
2
(13) Y(T−)=>5, Y(T−)=>6, X(T)=1,→Y(T)=5
In these rules, “A, B, C→D” means that if A, B and C hold, then D holds.
For example, the rule of (1) means that if the output before one time unit is 4 or less, the output before two time units is 5 or less, the current input is 0 and the input before one time unit is 0, then it is anticipated that the current output will become 6.
24
FIG. 14
It is now assumed to be requested to determine when what input should be given (input condition) in order to obtain Y=6 at a time later than the time in the time series data shown in .
FIG. 14
5
In the present embodiment, “logical inference” is performed by using the time series data shown in and the rules (1) to (13) in order to determine this input condition. This logical inference is performed by the condition decision section . Hereafter, this logical inference will be described.
FIG. 17
is a diagram showing the rules (1) to (13) in a table form.
FIG. 18
is a diagram explaining the logical inference.
FIG. 18
The logical inference predicts how the time series data changes after the next time while superposing at least the bottom end (last time) of the time series data on the rules as shown in .
FIG. 18
FIG. 14
23
2
23
2
24
1
25
In the example shown in , logical inference is performed by using the time series data shown in and the rule (9). In more detail, the value of Y in the time series data at time is 4, and the output at time T− in the rule (9) is “5 or less,” and consequently they match each other. Furthermore, the value of X in the time series data at time is 1, and the input at time T− in the rule (9) is 1, and consequently they match each other. In addition, the value of Y in the time series data at time is 5, and the output at time T− in the rule (9) is “5 or more,” and consequently they also match each other. If 0 is given as X at time (=T), therefore, it is anticipated that Y will become 5.
24
25
2
1
10
10
1
In the case of this example, the matched time zone (unified time zone) is two time units. In other words, the unified time zone includes time and in the time series data, and time T− and T− in the rule. As a matter of course, however, the unified time zone differs according to the size of the time zone included in the rule. If the time zones in the rule are T− to T, then, for example, ten time zones T− to T− are used.
24
FIG. 14
By using this logical inference, an input condition required to obtain Y=6 at a time later than the time in is determined.
FIG. 17
First, rules in which Y(T) is 6 are selected from among the rules (1) to (13) shown in , resulting in the rules (1), (4), (7) and (11).
FIG. 14
Subsequently, it is determined whether these rules (1), (4), (7) and (11) match the time series data shown in .
2
1
23
24
24
1
As for the rule (1), if time T− and T− in the rule (1) are respectively associated with time and in the time series data, then Y=5 at time does not satisfy Y<=4 at time T−. Therefore, the rule (1) does not match the time series data.
2
1
23
24
24
1
As for the rule (4), if time T− and T− in the rule (4) are respectively associated with time and in the time series data, then Y=5 at time does not satisfy Y<=4 at time T−. Therefore, the rule (4) does not match the time series data.
Determining for the rules (7) and (11) as well in the same way, neither of these rules matches the time series data.
Therefore, logical inference is performed by combining these rules.
In this case, rules are combined basically as a round robin. As a result, an input condition required to obtain Y=6 is determined by combining the rule (10) with the rule (4). A rule selection scheme to be used when combining rules is described in, for example, Journal of Information Processing Society of Japan, Vol. 25, No. 12, 1984.
FIG. 19
is a diagram showing concretely how logical inference is performed by combining the rule (10) with the rule (4).
2
1
1
2
1
23
24
FIG. 19
If time T− and time T− in the rule (4) are respectively associated with time T− and T in the rule (10) as shown in , then it will be appreciated that they match each other. Furthermore, if time T− and time T− in the rule (10) are respectively associated with time and in the time series data, then it will also be appreciated that they match each other.
25
26
If X=1 is given as the input at time , therefore, then it is anticipated that Y=4 will be outputted according to the rule (10). If X=1 is given as the input at time , then it is anticipated that Y=6 will be outputted according to the rule (4).
Processing steps performed by the inverse model calculation apparatus according to the present embodiment will now be described below.
FIG. 20
is a flow chart showing processing steps performed by the inverse model calculation apparatus.
2
1
71
First, the decision tree generation section generates a decision tree by using time series data recorded in the time series data recording section (step S).
2
5
72
Subsequently, the decision tree generation section gives an output value V at a future time (an output condition) to the condition decision section (step S).
5
73
74
The condition decision section empties the target list and the input list (step S), and adds an output condition “y(t)=V” to the target list as a target condition (step S).
5
75
The condition decision section executes a subroutine C described later (step S).
76
5
77
If a result of the execution of the subroutine C is false (YES at step S), then the condition decision section outputs a signal indicating that the given output value V cannot be obtained within a predetermined time (step S).
76
5
78
On the other hand, if the execution result of the subroutine C is true (NO at the step S), then the condition decision section outputs contents of the input list (input condition and value of time t) obtained in the subroutine C (step S).
FIG. 21
is a flow chart showing processing steps in the subroutine C in detail.
5
81
82
First, the condition decision section initializes a counter (for example, number of times i=0) (step S), and increments the i (i=i+1) (step S).
5
83
Subsequently, the condition decision section determines whether the number of times i has exceeded a predetermined value (step S).
83
5
84
If the i has exceeded the predetermined value (YES at the step S), then the condition decision section outputs a signal indicating that the given output value cannot be obtained (false) (step S).
83
5
85
On the other hand, if the i has not exceeded the predetermined value (NO at the step S), then the condition decision section determines whether a rule matching the time series data is present in the target list (step S).
5
85
86
At the current time, a rule is not stored in the target list. Therefore, the condition decision section determines that such an item is not present (NO at the step S), and takes out one item from the target list (step S).
5
87
The condition decision section determines whether the item taken out is an output condition or a rule (step S).
5
87
5
3
3
88
6
FIG. 16
If the condition decision section determines the item taken out to be an output condition (this holds true at the current time) (output condition at the step S), then the condition decision section causes the condition acquisition section to execute the subroutine A by using the item as the target condition, and receives a retrieval result (a rule including a value of the target condition in a leaf node) from the condition acquisition section (step S). For example, if the output value V is 5 in , then five rules (2), (5), (6), (9) and (13) are obtained by the subroutine A. If the output value V is , then four rules (1), (4), (7) and (11) are obtained.
89
5
84
If the retrieval result is false (YES at step S), then the condition decision section outputs a signal indicating that the given output value cannot be obtained (false) (step S).
89
5
3
90
On the other hand, if the retrieval result is not false (NO at the step S), then the condition decision section adds the rules acquired by the condition acquisition section to the target list (step S).
5
82
5
83
5
85
5
85
5
91
25
25
91
FIG. 17
FIG. 14
Subsequently, the condition decision section increments the i (step S). If the condition decision section determines that the i does not exceed a predetermined value (NO at the step S), then the condition decision section determines whether a rule that matches the time series data is present in the target list (step S). If the output value V is 5 in , then the rules (9) and (13) included in the rules (2), (5), (6), (9) and (13) match the time series data as shown in . In this case, the condition decision section determines that a matching rule is present (YES at the step ). The condition decision section specifies an input condition and time t on the basis of the matching rule and the time series data, and adds the input condition and the time t to the input list (step S). Here, X()=0 (rule (9)), X()=1 (rule (13)), and time t=25 are added to the input list (step S).
85
85
86
87
FIG. 17
On the other hand, if a rule matching the time series data is not present at the step S (NO at the step S), then one item is taken out from the target list (step S). For example, the rules (1), (4), (7) and (11) in the case where the output value V is 6 in do not match the time series data. Therefore, one of these items (rules) is taken out from the target list. Here, for example, the rule (4) is taken out (rule at the step S).
5
3
92
The condition decision section causes the condition acquisition section to determine whether a rule that matches the rule taken out (object rule) is present (step S).
92
5
93
FIG. 17
If such a rule is present (YES at the step S), then the condition decision section adds that rule to a temporary list together with the above-described object rule (step S). If the output value V is 6 in , then rules (10) and (13) are present as rules matching the rule (4). Therefore, the rule (4) serving as the object rule, and the rules (10) and (13) obtained as matching the rule (4) are stored in the temporary list.
5
94
5
The condition decision section determines whether the obtained rules in the temporary list match the time series data (step S). In the above described example, the condition decision section determines whether the rule (10) or the rule (13) matches the time series data.
94
5
96
5
25
26
If a matching rule is present (YES at step S), then the condition decision section specifies the input condition and the time t on the basis of the matching rule and the object rule, and adds the input condition and the time t to the input list (step S). For example, in the above-described example, the condition decision section specifies X()=1 as the input condition on the basis of the rule (10) and X()=1 as the input condition on the basis of the rule (4), and adds these input conditions to the input list together with time t=26.
5
97
97
5
97
5
82
The condition decision section determines whether the target list is empty (step S). If the target list is empty (YES at the step S), then the condition decision section terminates the subroutine C. If the target list is not empty (NO at the step S), then the condition decision section empties the temporary list, and returns to the step S.
94
94
5
92
93
92
5
93
92
5
95
82
If the obtained rule in the temporary list does not match the time series data at the step S (NO at the step S), then the condition decision section performs the steps S and S again by using the rule that does not match as an object rule. If a rule that matches the object rule is obtained (YES at the step S), then the condition decision section adds the rule to the temporary list (step S). On the other hand, if a rule is not obtained (NO at the step S), then the condition decision section empties the temporary list (step S), and returns to the step S.
According to the present embodiment, a condition required to obtain a given output value is calculated by combining rules obtained from the decision tree so as to go back in time. Therefore, condition calculation can be terminated in a short time.
(Fifth Embodiment)
2
1
In the fourth embodiment, the whole time zone except the current time T is used as the time zone of matching between rules and matching between a rule and time series data, i.e., the time zone of unification. In the fourth embodiment, the time zone of the unification is two time units ranging from T− to T−. If rules are unified in the whole time zone except the current time in the case where the time zone included in the rules is long, then a high-precision inference is anticipated, but a large amount of calculation is requested, resulting in inefficiency in many cases. If unification can be performed in a shorter time zone, then the efficiency is high. If the time zone of unification is made shorter, however, a problem to which inference precision fall may occur. In the present embodiment, therefore, a value effective as a time zone of unification is calculated and unification is performed with that value, and thereby inference is implemented with a small amount of calculation and with high precision.
First, the relation between the time zone of unification and the inference precision will be described briefly.
1
2
1
1
2
1
The relation will now be described by taking the rule (4) as an example. As described above, “Y(T−)<4, Y(T−)<=5, X(T)=1, X(T−)=1→Y(T)=6” in the rule (4) means that the result on the right side (the value of the object variable) is obtained when all conditions (conditions of the explaining variables) on the left side in this logical expression hold. If X(T−)=1 is set after Y(T−)<=5 has held, then it is indistinct from the rule (4) whether Y(T−)<=4. In other words, it is indistinct whether the value of Y at each time in the rule holds if a condition before that time and at that time has held.
In the present embodiment, a probability (stochastic quantity) that an output condition at each time included in the rule will hold in the case where conditions before that time and at that time hold is found, and unification is performed in a minimum time zone having the probability higher than a threshold. As a result, it can be anticipated to perform logical inference with a minimum calculation quantity and high precision. Hereafter, this will be described in more detail by taking the rule (4) as an example.
FIG. 14
Hereafter, the probability that an output condition at each time included in the rule (4) will hold in the case where conditions before that time and at that time hold will be described by using the time series data shown in .
2
First, as for Y(T−)<=5 in the rule (4), other conditions before this time and at this time are not present, and consequently it will be omitted.
1
1
2
4
13
19
23
10
14
18
20
22
1
FIG. 14
Subsequently, as for Y(T−)<=4, it is checked whether it holds assuming that X(T−)=1 when Y(T−)<=5 holds. As a result, it holds at time , , and , and it does not hold at time , , , and in the time series data in . Therefore, the probability that Y(T−)<=4 will hold is 44% (=4/9×100%).
2
1
Therefore, as for the rule (4), if the threshold is set equal to 40%, it can be the that unification using two time zones (T−, T−) is suitable.
FIG. 22
FIG. 21
89
Processing steps of calculating time zones in which unification is performed and performing the unification in the calculated time zones will now be described. This is achieved by executing a subroutine D shown in instead of the step S shown in .
FIG. 22
is a flow chart showing processing steps in the subroutine D.
3
101
5
3
1
102
5
102
5
90
85
92
94
5
92
5
FIG. 21
If a result of retrieval performed by the condition acquisition section is not false (NO at step ), then the condition decision section calculates the probability that an output condition at each time in each of the rules acquired from the condition acquisition section will hold when a condition at an earlier time and at the time holds, on the basis of the time series data in the time series data recording section (step S). The condition decision section sets a minimum time zone having the probability greater than a threshold as the time zone for unification (step S). The condition decision section adds each retrieved rule to the target list together with the time zone of unification of each rule (step S). In the steps S, S and S of performing unification (see ), the condition decision section performs unification by using the calculated time zone. If a new rule is acquired at the step S, then the condition decision section finds a time zone in the same way.
3
101
3
84
On the other hand, if the result of the retrieval performed by the condition acquisition section is false (YES at the step S), then the condition decision section proceeds to the step S, and outputs a signal (false) indicating that the given output value V cannot be obtained.
102
3
At the above-described step S, the time zone of unification has been calculated for each of rules. However, a time zone common to all rules may be found. Specifically, the condition decision section calculates an average of holding probability of the output condition at each time for all rules, and uses a time zone having the average exceeding the threshold as a time zone common to the rules.
FIG. 23
FIG. 21
81
82
This is implemented by adding a subroutine E shown in between, for example, steps S and S shown in .
5
3
5
5
112
5
85
92
94
FIG. 21
In other words, the condition decision section causes the condition acquisition section to acquire all rules included in the decision tree. The condition decision section calculates the holding probability of the output condition at each time with respect to all acquired rules, and finds an average of the holding probability at each time. The condition decision section specifies time when the value becomes equal to the threshold or more, and sets a time zone before the specified time (including the specified time) as the time zone of unification common to the rules (step S). Therefore, the condition decision section uses this common time zone at the steps S, S and S shown in .
According to the present embodiment, a minimum time zone satisfying predetermined precision is adopted as the time zone for unification, as heretofore described. Therefore, the processing can be executed by using a small quantity of calculation without lowering the precision much. Furthermore, according to the present embodiment, a time zone for unification common to the rules is calculated. Therefore, the processing efficiency can be further increased.
(Sixth Embodiment)
In fields of control or the like, there are a plurality of process outputs in many cases. There is a case where it is desirable to perform inverse calculation for a plurality of outputs. In other words, there is a case where it is desirable to find an input that makes a plurality of outputs simultaneously desirable values, for example, an input that makes the temperature of an apparatus and the pressure of another apparatus connected to the apparatus simultaneously desirable values.
As a first method, there is a method of converting a plurality of outputs to a one-dimensional evaluation value and constructing a model for the one-dimensional evaluation value. In the case where the evaluation value is one-dimensional, it is possible to construct a decision tree and execute inverse calculation by using the constructed decision tree.
In this method, however, a proper evaluation function for conversion to a one-dimensional evaluation value must be defined. The proper evaluation function differs depending upon the problem, and it is difficult to properly define the evaluation function. Even if an evaluation function can be defined properly, since conversion processing to the evaluation value exists for constructing a model, this method results in a problem of a prolonged calculation time.
As a second method, a method of regarding a direct product (set) of a plurality of outputs as a value of one object variable and constructing a model such as a decision tree is conceivable.
If in this method a loss (blank) is present in the value of the object variable in the observed data, then data of that portion cannot be used for construction of the decision tree. In other words, only data having complete values of all object variables can be used for constructing decision tree. Therefore, in this method, there is a fear that usable data will be remarkably limited. Fewer data used for construction exert a bad influence upon the precision of the generated decision tree, and there is also a fear that the decision tree will not be useful.
As a third method, there is a method of generating a plurality of decision trees with respect to each of a plurality of outputs and performing inverse calculation by using a plurality of decision trees simultaneously.
However, this method is difficult, or requires a long calculation time. The reason can be explained as follows. Even if a value of an explaining variable that makes certain one object variable desirable value is found by using certain one decision tree, the value of the explaining variable does not always satisfy the condition with respect to a different object variable.
In view of the problems heretofore described, the present inventors have gone through unique studies. As a result, the present inventors have acquired a technique of combining decision trees generated for respective object variables and generating a composite decision tree having a set of these object variables as an object variable. In other words, this composite decision tree has, in its leaf node, a value obtained by combining values of leaf nodes in decision trees. A condition required to simultaneously obtain a plurality of desirable outputs can be calculated by applying this composite decision tree to the first to fifth embodiments. Hereafter, the technique for combining the decision trees will be described in detail.
FIG. 25
is a configuration diagram of a decision tree combination apparatus, which combines a plurality of decision trees.
11
12
13
14
The decision tree combination apparatus includes a data input section , a decision tree generation section , a decision tree combination section , and a decision tree output section .
11
12
FIG. 2
The data input section inputs data including a value of an explaining variable and values of object variables to the decision tree generation section . The value of the explaining variable is, for example, an operation value inputted into a device. The values of the object variables are resultant outputs (such as the temperature and pressure) of the device. The present data includes a plurality of kinds of object variables. Typically, the data are collected by observation and recording (see ).
12
12
12
The decision tree generation section generates one decision tree on the basis of the value of the explaining variable included in the data and the value of one of object variables included in the data. The decision tree generation section generates one decision tee for each of the object variables in the same way. In other words, the decision tree generation section generates as many decision trees as the number of the object variables. Each decision tree has a value of an object variable at a leaf node (terminal node). Nodes other than leaf nodes become explaining variables. A branch that couples nodes becomes a value of an explaining variable.
13
12
1
2
3
1
2
3
1
1
1
2
1
2
1
2
2
2
1
2
1
2
1
The decision tree combination section combines a plurality of decision trees generated in the decision tree generation section , and generates one decision tree (composite decision tree) that simultaneously infers values of a plurality of object variables on the basis of the value of the explaining variable. This composite decision tree has, at its leaf node, a set of values of object variables obtained by combining values of leaf nodes (values of object variables) in the decision trees. For example, assuming that a first decision tree has y, y, y, . . . yn at respective leaf nodes and a second decision tree has z, z, z, . . . zn at respective leaf nodes, leaf nodes of the combined decision tree become (y,z), (y,z) . . . (y,zn), (y,z), (y,z), . . . (yn,zn). By using this composite decision tree as the object decision tree in the above-described first to fifth embodiments, a condition required to satisfy the values of a plurality of object variables simultaneously can be found. For example, when using this composite decision tree in the first embodiment and obtaining (y,z) as an output value at a given future, a condition required to obtain this value (y,z) can be found by specifying a leaf node having the value (y,z) and tracing branches from this leaf node toward the root node.
14
13
3
FIGS. 1 and 9
The decision tree output section outputs the composite decision tree generated by the decision tree composite section . The outputted composite decision tree can be used as the object decision tree in the first to fifth embodiments. In other words, the condition acquisition section shown in can use this composite decision tree as the object decision tree.
FIG. 25
Hereafter, the apparatus shown in will be described in detail by using a concrete example.
FIG. 27
is a table showing an example of observed data.
1
2
3
4
5
6
1
2
1
2
3
4
5
6
1
2
1
6
1
2
1
6
1
2
There are a large number of instances, such as an instance having 1 as the value of variable X, 2 as the value of variable X, 0 as the value of variable X, 0 as the value of variable X, 0 as the value of variable X, A as the value of variable X, 7 as the value of variable Y and A as the value of variable Y, and an instance having 3 as the value of variable X, 0 as the value of variable X, 1 as the value of variable X, 0 as the value of variable X, 1 as the value of variable X, B as the value of variable X, 7 as the value of variable Y and C as the value of variable Y. Here, X to X are explaining variables, and Y and Y are object variables. In the field of control, values of X to X correspond to the input into a target system (such as an item representing the material property and operation value of the device), and values of Y and Y. correspond to the output from the target system (such as the temperature and pressure of a material).
FIG. 27
11
12
First, data shown in are inputted from the data input section to the decision tree generation section . The inputted data are stored in a suitable form.
12
Subsequently, in the decision tree generation section , a decision tree is generated per object variable.
11
1
FIG. 27
FIG. 28
If data inputted from from the data input section are the data shown in , then there are two object variables, and consequently two decision trees are generated. Data used to generate one decision tree (a decision tree associated with the object variable Y) are shown in .
FIG. 28
FIG. 27
2
1
The data shown in are obtained by deleting the data of the object variable Y and leaving the data of the object variable Y in the data shown in .
A method used to generate a decision tree on the basis of data thus including only one object variable is described in, for example, “Data analysis using AI” written by J. R. Quinian, translated by Yasukazu Furukawa, and published by Toppan Corporation in 1995, and “Applied binary tree analysis method” written by Atsushi Otaki, Yuji Horie and D.Steinberg and published by Nikks Giren in 1998.
2
1
FIG. 27
In the same way, the decision tree associated with the object variable Y can also be generated. Data used to generate this decision tree are obtained by deleting the data of the object variable Y in the data shown in .
1
2
1
2
Decision trees generated for the object variables Y and Y as heretofore described are herein referred to as “decision tree ” and “decision tree ” for convenience.
FIG. 26
12
12
12
12
12
a
b
a
b
Here, as shown in , which shows another example of the decision tree combination apparatus, it is also possible to divide the decision tree generation section into a data shaping processing section and a decision tree generation processing section , cause the data shaping processing section to generate data including only one object variable, and cause the decision tree generation processing section to generate a decision tree by using the data. Decision trees associated with object variables may be generated in order or may be generated in parallel.
FIG. 28
Although in generating a decision tree of each object variable, data including only one object variable has been generated temporarily (see ), this processing is performed in order to simplify the description and therefore it may be omitted in the actual processing.
FIG. 29
1
2
1
2
is a diagram showing examples of the decision tree and the decision tree generated for the object variables Y and Y.
1
2
Hereafter, how to see the decision tree and the decision tree will be explained briefly.
1
1
1
1
3
3
1
3
1
1
2
5
FIG. 29
The decision tree classifies the instance according to the value of Y, which is an object variable (leaf node). First, it is determined whether X is greater than 4. If X is equal to 4 or less, then it is determined whether X is 0 or 1. If X is equal to 0, then Y is determined to be less than 2. If X is equal to 1, then Y is determined to be greater than 5. Also when X is greater than 4, similar processing performed. In , “2-5” in a leaf node means “between 2 and 5 inclusive of and .”
2
2
3
3
4
4
2
4
2
3
In the same way, the decision tree classifies the instance according to the value of Y. First, it is determined whether X is 0 or 1. If X is 0, then it is determined whether X is 0 or 1. If X is 0, then Y is determined to be A. If X is 1, then Y is determined to be C. Also when X is 1, similar processing is performed.
1
2
1
2
FIG. 27
These decision trees and classify instance sets included in already known data (see ). Even for new data, however, values of Y and Y, which are object variables, can be predicted.
Typically, classification using a decision tree is not right a hundred percent. Because there is a contradiction in data used to construct a decision tree, in some cases. Furthermore, because an instance that occurs only a few times is regarded as an error or noise and it does not exert an influence upon the construction of a decision tree in some cases. It is possible to generate a detailed decision tree that correctly classifies data obtained at the current time hundred percent, but actually such a decision tree is not so useful. Because it is considered that such a decision tree faithfully represents even noise and errors. In addition, such a decision tree merely re-represents the current data strictly, and necessity of re-representing the current data in a decision tree form is weak. Furthermore, because a decision tree that is too detailed becomes hard for the user to understand. Therefore, it is desirable to generate a compact decision tree with processing performed for noise moderately.
13
The decision tree combination section combines a plurality of decision trees as described above and generates one decision tree. Hereafter, three kinds (combination methods 1 to 3) of concrete example of decision tree combination method will be described. However, it is also possible to use a combination of them.
Hereafter, the combination methods 1 to 3 will be described in order.
(Combination Method 1)
FIG. 30
is a flow chart showing a processing procedure for performing the combination method 1.
1001
1
2
3
4
5
6
FIG. 27
FIG. 31
In the combination method 1, first, a series of values of explaining variables (explaining variable values) is generated (step S). The “series of explaining variable values” means, for example, input data having values of the explaining variables X, X, X, X, X and X shown in . First, one series is generated. It is now assumed that a series of explaining variable values shown in has been generated.
1
2
1002
1003
Subsequently, the decision trees and are provided with the series of explaining variable values, and the value of the object variable is obtained (steps S and S). In other words, a certain leaf node is arrived at by tracing a decision tree from its root node in order. The value of the leaf node is the value of the object variable.
1
1
1
3
2
3
4
Specifically, in the decision tree , X is 1, i.e., X is “<=4,” and consequently the processing proceeds to a left-side branch. Subsequently, since X is 0, the processing proceeds to a left-side branch. As a result, a leaf node of “<2” is arrived at. On the other hand, in the decision tree , X is 0, and consequently the processing proceeds to a left-side branch. Subsequently, since X is 0, the processing proceeds to a left-side branch. As a result, a leaf node of “A” is arrived at.
1
2
1004
FIG. 31
FIG. 32
The values of the leaf nodes thus obtained from the decision trees and are added to the table shown in to generate one instance (step S). shows one generated instance data.
1
2
Subsequently, a different series of explaining variable values is generated. In this case as well, there are no constrain in how to generate the series, but it is desirable that the generated series is not the same as the series generated earlier. It is desirable to generate all combinations of explaining variable values by changing the values of explaining variables, for example, at random or in order. The series generated is given to the decision trees and to acquire the values of the object variables and obtain instance data. By repeating the above, a set of instance data is generated.
1005
FIG. 32
A decision tree is generated by using the set of generated instance data and regarding a set of two object variables as one object variable (step S). For example, a decision tree is generated by regarding “<2” and “A” in as the value of one object variable. Since the decision tree generation method is described in the above-described document, here detailed description will be omitted.
(Combination Method 2)
FIG. 33
is a flow chart showing a processing procedure for performing a combination method 2.
1
2
1011
1012
First, paths (rules) from the root node to leaf nodes are acquired from each of the decision trees and , and all combinations of the acquired paths are generated. As a result, a plurality of path sets are generated. And by, for example, concatenating paths included in each path set, one new path (composite path) is generated from each path set, thereby, a new path set (a set of composite path) is obtained (step Subsequently, composite paths included in the new path set obtained at the step S is combined to obtain one decision tree (step S).
1011
1012
Hereafter, the steps S and S will be described in more detail.
1011
First, the step S will be described.
FIG. 34
1011
is a flow chart showing a processing procedure at the step S.
1
2
1
2
1021
First, paths from the root node to leaf nodes are acquired from each of the decision trees and . The acquired paths are combined between the decision trees and in every kinds of combination, and a plurality of path sets are generated (step S).
FIG. 35
FIG. 35
FIG. 29
FIG. 35
1
2
is a diagram showing an example of a path set. The left side of shows a path from the root node of the decision tree (see ) to the leftmost leaf node, and the right side of shows a path from the root node of the decision tree to the leftmost leaf node. Each path does not include branching.
1
2
1
2
Paths included in the decision tree and paths included in the decision tree are thus combined successively. It does not matter which order paths are combined in. However, all combinations are performed. Since the decision tree has five leaf nodes and the decision tree has six leaf nodes, (5×6=) 30 path sets are obtained.
1022
FIG. 34
Upon thus acquiring path sets, paths included in each path set are concatenated longitudinally to generate a new path (concatenated path) (step S in ).
FIG. 36
FIG. 35
is a diagram showing a state in which the path set shown in have been concatenated.
FIG. 36
2
1
1
2
The leaf nodes (object variables) in paths before concatenation are assigned to an end of the concatenated path. Other nodes (explaining variables) are concatenated in the longitudinal direction. In , the path of the decision tree is concatenated under the path of the decision tree . However, the path of the decision tree may also be concatenated under the path of the decision tree .
1023
FIG. 34
Subsequently, it is checked whether there is a contradiction in the concatenated path (step in ).
The “contradiction” means that there are duplicating explaining variables and their values are different from each other. For example, if two or more same explaining variables (nodes) are included in the concatenated path and one of them is 1 whereas the other is 0, then there is a contradiction.
1023
1024
1026
3
3
FIG. 36
If there is a contradiction (YES at the step S), then this concatenated path is deleted (step S), and the next path set is selected (YES at step S). In , there are two nodes X. Since the two nodes X have the same value 0, there is no contradiction.
1023
1025
1023
3
3
1025
FIG. 36
FIG. 36
FIG. 37
If there is no contradiction (NO at the step S), then processing for eliminating duplication included in the concatenated path is performed (step S). The “duplication” means that there are a plurality of same explaining variables (nodes) in the concatenated path and the explaining variables have the same value. The contradiction check has been performed at the step S. If there are a plurality of same explaining variables at the current time, therefore, the explaining variables should have the same value, and consequently there is duplication. If there is duplication, a duplicating explaining variable (node) and its branch are deleted from the concatenated path. As a result, the concatenated path becomes shorter. In , two nodes X are included in the concatenated path, and the two nodes X have a value of 0. Therefore, this is duplication. A path (composite path) obtained by eliminating the duplication from the concatenated path shown in is shown in . The path generated by the step S is referred to as “composite path”.
1022
1024
1025
1024
FIG. 38
As heretofore described, the concatenation processing (the step S), the contradiction processing (the step S), and the duplication processing (the step S) are performed for each path set (30 path sets in the present example). Since contradicting concatenated paths are deleted by the contradiction processing (the step S), the number of generated composite paths becomes 30 or less. In the present example, 16 composite paths are generated. shows 16 generated composite paths.
FIG. 38
FIG. 38
FIG. 38
1
2
1
2
1
2
1
3
1
4
1024
1
2
In , parenthesized numerical values indicated above each composite path represents how paths in the decision tree and the decision tree have been combined. For example, (-) means that a path including the leftmost leaf node in the decision tree and a path including the second leftmost leaf node in the decision tree have been combined. In , (-) and (-) are not present, because they have been deleted by the above-described contradiction processing (step S). In each composite path, nodes may be interchanged in arrangement order except the leaf nodes (object variables). In , nodes are arranged in the order of increasing number like X, X, . . . , for easiness to see.
1024
1025
Furthermore, the contradiction processing (the step S) and the duplication processing (the step S) may be inverted in execution order, or they may be executed in parallel. In this case as well, the same result is obtained.
1012
FIG. 33
The step S (see ) will now be described in detail.
1012
FIG. 38
At the step S, one decision tree is constructed by combining the composite paths (see ) generated as heretofore described.
FIG. 39
1012
is a flow chart showing the processing procedure at the step S in detail.
1031
FIG. 38
First, all composite paths are handled as objects (step S). In the present example, 16 composite paths shown in are handled as objects.
1032
Subsequently, it is determined whether there are two or more object composite paths (step S). Since there are 16 object composite paths at the current time, the processing proceeds to “YES.”
1033
1
3
1
1
2
1
3
1
2
FIG. 38
Subsequently, an explaining variable (node) that is included most among the set of object composite paths is selected (step S). Upon checking 16 composite paths, it is found that the nodes X and X are used in all composite paths, and included most (respectively 16 times). If there are a plurality of most nodes, then arbitrary one node is selected. It is now assumed that the node X is selected. By the way, the composite paths shown in are generated on the basis of the decision tree and the decision tree . Therefore, each composite path necessarily includes the root nodes (the nodes X and X in the present example) of the decision trees and .
1034
1
Subsequently, the selected node is coupled under a branch selected in a new decision tree (the decision tree in the middle of generation), as a node of the new decision tree (step S). In first processing (a loop of the first time), however, the node is designated as a root node. At the current time, therefore, the node X is designated as the root node.
1035
1
1
FIG. 38
FIG. 40
Branches are generated for the node on the basis of values that the node can have (step S). The values that the node can have are checked on the basis of a set of composite paths. Checking the values that the node X can have on the basis of the set of composite paths shown in , “<=4” and “4<” are obtained. Therefore, branches of “<=4” and “4<” are generated for the node X. The decision tree in the middle of generation generated by processing heretofore described is shown in .
1036
FIG. 40
Subsequently, one branch is selected in the decision tree at the current time (step S). It is now assumed that the left-hand “<=4” branch has been selected in . The right-hand branch will be subject to processing later. Either branch may be selected earlier.
FIG. 38
FIG. 38
1036
1037
1
1
Subsequently, the set of composite paths shown in is searched for composite paths including a path from the root node of this decision tree to the branch selected at the step , and found paths are designated as object composite paths (step S). In the present example, composite paths including “X<=4” are searched for, and the composite paths are designated as object composite paths. In the set of composite paths shown in , composite paths including “X<=4” are six composite paths shown in the highest column. Therefore, these six composite paths are designated as object composite paths.
1032
Returning back to the step S, it is determined whether there are two or more object composite paths. Since there are six object composite paths, the processing proceeds to “YES.”
1033
1037
1
1036
3
3
FIG. 38
Subsequently, a node that is included most among the set of object composite paths is selected (step S). Here, however, the node used to search for object composite paths at the step S (the node X in the present example), i.e., the node on the path from the root node of the decision tree to the branch selected at the step S is excluded. Since a node that is most included among the six composite paths shown in the highest column of is X, the node X is selected.
1036
1034
1036
3
FIG. 40
Subsequently, the selected node is coupled under the branch selected at the step S, as a node of the new decision tree (step S). Since the branch selected at the step S is the left-hand branch shown in , the node X is coupled under the branch.
1035
3
3
FIG. 41
Branches are generated for the node on the basis of values that the coupled node can have (step S). Since the values that the node X can have are “0” and “1,” branches of “0” and “1” are generated under the node X. The decision tree generated heretofore is shown in .
1036
3
Subsequently, one branch is selected in the decision tree (step S). It is now assumed that the left-hand “0” branch has been selected from branches branched from the node X.
1036
1037
1036
3
1
3
FIG. 38
Subsequently, the set of composite paths (six composite paths shown in the highest column) is searched for composite paths including a path from the root node of this decision tree to the branch selected at the step , and found paths are designated as object composite paths (step S). The branch selected at the step S is the left-hand “0” branch in branches branched from the node X. Therefore, the six composite paths shown in the highest column are searched for composite paths including paths (“X<=4” and “X=0”) from the root node to that branch. Two composite paths, i.e., the leftmost composite path and the second leftmost composite path shown in the highest column of are paths satisfying the above condition.
1032
Returning back to the step S, it is determined whether there are two or more object composite paths. Since there are two object composite paths, the processing proceeds to “YES.”
1033
1
3
1
3
4
4
Subsequently, a node that is included most among the set of object composite paths is selected (step S). However, the nodes X and X are excluded. Excluding the nodes X and X, node included most in two object composite paths is the node X, and consequently the node X is selected.
1036
1034
1036
3
4
3
FIG. 41
Subsequently, the selected node is coupled under the branch selected at the step S, as a node of the new decision tree (step S). Since the branch selected at the step S is the left-hand branch (X=0) shown in , the node X is coupled under the “0” branch branched from the node X.
1035
4
4
FIG. 38
FIG. 42
Branches are generated for the node on the basis of values that the coupled node can have (step S). The values that the node X can have are “0” and “1” respectively on the basis of the leftmost composite path and the second leftmost composite path shown in the highest column of . Therefore, branches corresponding to “0” and “1” are generated under the node X (see ).
1036
4
Subsequently, one branch is selected in the decision tree (step S). It is now assumed that the left-hand “0” branch has been selected from branches branched from the node X.
FIG. 38
FIG. 38
1036
1037
Subsequently, the set of composite paths shown in is searched for composite paths including a path from the root node of this decision tree to the branch selected at the step , and found paths are designated as object composite paths (step S). The composite path that becomes the object is only the leftmost composite path in the highest column shown in .
1032
Returning back to the step S, it is determined whether there are two or more object composite paths. Since there is only one object composite path, the processing proceeds to “NO.”
1036
1038
FIG. 42
Subsequently, the leaf node in this composite path is coupled under the branch selected at the step S, and designated as a leaf node of the new decision tree (step S). In the present example, “<2, A” becomes the leaf node of the new decision tree. The decision tree generated heretofore is shown in .
1039
FIG. 42
Subsequently, it is determined whether there is a branch that is not provided with a leaf node in the decision tree (step S). Since there are three branches having no leaf nodes as shown in , the processing proceeds to “YES.”
1040
4
FIG. 42
Subsequently, one branch having no leaf node is selected in this decision tree (step S). It is now assumed that a branch of “X=1” has been selected in the decision tree shown in . The selected branch may be any branch so long as it has no leaf node.
1037
1040
FIG. 38
FIG. 38
Subsequently, the processing proceeds to the step S. The set of composite paths shown in is searched for a composite path including a path from the root node to the branch selected at the step S in the decision tree at the current time, and the found composite path is designated as the object composite path. Here, only the second leftmost composite path in the highest column of is designated as the object composite path.
1032
Returning back to the step S, it is determined whether there are two or more object composite paths. Since there is only one object composite path, the processing proceeds to “NO.”
1040
FIG. 43
Subsequently, a leaf node in this composite path is coupled under the branch selected at the step , and it is designated as a leaf node in the new decision tree. In the present example, “<2, C” becomes a leaf node in the new decision tree. The decision tree generated heretofore is shown in .
1
2
FIG. 44
By continuing similar processing thereafter, a decision tree obtained by combining the decision tree with the decision tree is finally generated as shown in .
1033
FIG. 39
With reference to the step S shown in , it has been described that any node may be selected if there are nodes having the same number, when finding a node that is included most in the set of object composite paths. There may be a doubt that the finally obtained decision trees may be different. However, finally obtained decision trees become equal in meaning. Because even if such a node is not selected at a certain time, the node is certainly selected at the next or subsequent selection opportunity. Since a leaf node of the new decision tree is generated on the basis of a combination of leaf nodes of both decision trees, contents of the finally obtained decision tree do not depend upon the order of node selection.
(Combination Method 3)
FIG. 45
is a flow chart showing a processing procedure for performing a combination method 3.
1041
1
2
1
3
FIG. 29
First, as represented by a step S, root nodes respectively of the decision tree and the decision tree are handled as objects. In the present example, the nodes X and X become objects (see ).
1042
1
3
1
3
1
3
1
2
1
3
1
3
Subsequently, object nodes are combined between different decision trees to generate a node set. The set of nodes are designated as node of a new decision tree (step S). In the present example, the set of the nodes X and X is designated as a node (set node) of the new decision tree. This node is referred to as “X, X”. Unless this set node is composed of a leaf node, a node corresponding to this set node is detected from each decision tree, and branches of detected nodes are combined to generate a new branch. The generated new branch is added to the set node. In the present example, nodes corresponding to the node “X, X” in the decision tree and the decision tree are X and X. Therefore, branches of the nodes X and X are combined to generate a new branch.
1
1
3
2
1
3
FIG. 46
For more detail, the node X in the decision tree has branches of “<=4” and “4<”, and the node X in the decision tree has branches of “0” and “1.” Therefore, four new branches of “<=4, 0”, “<=4, 1”, “4<, 0” and “4<, 1” are generated and added to the node “X, X.” The decision tree in the middle of generation generated heretofore is shown in .
1043
FIG. 46
Subsequently, it is determined whether there is a branch having no leaf node (step S). As shown in , there are four branches having no leaf node, and consequently the processing proceeds to “YES.”
1044
Subsequently, one branch having no leaf node is selected (step S). It is now assumed that the leftmost branch has been selected. However, the selected branch may be any branch.
1
2
1045
1
3
1
1
1
3
3
3
2
1
3
4
FIG. 46
Subsequently, a branch of the decision tree and a branch of the decision tree corresponding to the selected branch are detected, and a node following this branch is selected as an object (step S). As described above, the selected branch is the leftmost branch shown in , i.e., a branch of “X<=4, X=0.” Therefore, a branch “X<=4” in the decision tree corresponding to the branch of “X<=4, X=0” is traced, and the next node X is selected. In the same way, a branch “X=0” in the decision tree corresponding to the branch of “X<=4, X=0” is traced, and the next node X is selected. These nodes are designated as objects.
1042
3
4
3
4
3
4
1042
3
4
FIG. 46
FIG. 47
Returning back to the step S, nodes designated as the objects are combined to generate a new node. This new node is added to the new decision tree. In the present example, the nodes designated as the objects are X and X. In , therefore, a node “X, X” is added under the leftmost branch. In the same way as the foregoing description, branches are branched from that node “X, X”. As a result, branches of four kinds, i.e., branches of “0, 0”, “0, 1”, “1, 0” and “1, 1” are added (step S). The decision tree heretofore generated is shown in . Due to the restriction imposed on paper, only the leftmost branch, among branches branched from the node “X, X” is provided with values.
1043
Subsequently, it is determined whether there is a branch having no leaf node in the decision tree at the current time (step S). Since any branch is not yet provided with a leaf node, the processing proceeds to “YES.”
1044
Subsequently, one branch having no leaf node is selected (step S). It is now assumed that the leftmost branch has been selected.
1
2
1045
3
1
4
2
FIG. 47
FIG. 47
FIG. 47
Subsequently, a branch of the decision tree and a branch of the decision tree corresponding to the selected branch are specified, and a node following this branch is selected as the object (step S). In the present example, the leftmost branch in has been selected. Therefore, a node “<2” following a branch “X=0” in the decision tree corresponding to the leftmost branch in , and a node “A” following a branch “X=0” in the decision tree corresponding to the leftmost branch in are selected.
1042
1042
1
2
Returning back to the step S, nodes designated as the objects are combined to generate a new node. This new node is added to the new decision tree (step S). In the present example, a node “<2, A” is added as a new node. Since the nodes “<2” and “A” are leaf nodes in the decision tree and the decision tree , however, the newly generated node “<2, A” becomes a leaf node in the new decision tree. Therefore, branched branches are not generated from the node “<2, A.” If at this time one of the nodes is a leaf node in the original decision tree, whereas the other of the nodes is not a leaf node, then branched branches are further generated by using the decision tree including the node that is not a leaf node, in the same way as the foregoing description.
FIG. 48
By repeating the processing heretofore described, a decision tree shown in is finally generated.
FIG. 48
FIG. 48
In , parts of the tree are enlarged and shown in different places due to the restriction on the paper space. In , a path provided with a mark “X” is not actually present because there is a contradiction, but it is shown in order to express clearly the fact.
Heretofore, the combination methods 1, 2 and 3 have been described. The combination method 2 and the combination method 3 produce decision trees that are equal in meaning. It is a possibility that the combination method 1 will produce a decision tree that is slightly different from that produced by the combination method 2 and the combination method 3, depending upon given data. If the number of data is large, however, there is no great difference.
An improvement method for the decision tree generated as heretofore described will now be described below.
Typically, a decision tree has not only information concerning the branches and nodes, but also various data calculated to construct the decision tree from observed data. Specifically, the decision tree has the number of instances in each explaining variable (node) (for example, when a certain explaining variable can has “0” and “1” as its value, the number of instances in the case of “0” and the number of instances in the case of “1”), and distribution of the number of instances in each explaining variable with respect to the value of an object variable (for example, when there are 100 instances in which a certain explaining variable becomes “0” in value, there are 40 instances in which the object variable becomes A in value and 60 instances in which the object variable becomes B in value). By using these kinds of information hold by the decision tree, therefore, a composite decision tree generated by using one of the combination methods 1 to 3 is evaluated, and the composite decision tree is improved by deleting paths having low precision.
FIG. 49
FIG. 48
FIG. 4
1
2
is a diagram showing an evaluation method of a leftmost path in the composite decision tree (see ). The leftmost path in is a path generated by combining leftmost paths respectively in the decision tree and the decision tree .
FIG. 49
1
1
3
1
The left side of shows the leftmost path of the decision tree . There are 100 instances satisfying “X<=4” and “X=0”. There are 70 instances in which the value of the object variable becomes “<2”, 20 instances in which the value of the object variable becomes “2-5” (between 2 and 5 inclusive 2 and 5), and 10 instances in which the value of the object variable becomes “5<”. In other words, the precision of the path in the decision tree is 70% (70/100).
FIG. 49
2
3
4
2
The right side of shows the leftmost path of the decision tree . There are 90 instances satisfying “X=0” and “X=0”. There are 80 instances in which the value of the object variable becomes “A”, and 20 instances in which the value of the object variable becomes “B”. In other words, the precision of the path in the decision tree is 80% (80/100).
1
3
4
When “X<=4” and “X=0” and “X=0”, therefore, it is inferred that the probability of the value of the object variable becoming “<2, A” is 70%×80%=56%.
1
2
FIG. 49
By the way, it is impossible that the number of instances in the composite decision tree becomes greater than the number of instances in the original decision tree. Therefore, the number of instances in the composite decision tree becomes at most, min{the number of instances in the composite decision tree , the number of instances in the composite decision tree }. In the present example, the number of instances in the composite decision tree becomes 90 or less as shown in .
1
3
4
On the basis of this, in the composite decision tree, when “X<=4” and “X=0” and “X=0”, it is inferred that the number of instances in which the value of the object variable becomes “<2, A” is at most 90×56%=approximately 50. If this value or probability is equal to a predetermined value or less, then the composite decision tree is improved by deleting.
Furthermore, it is also possible to apply each path (rule corresponding to each path) of the composite decision tree to already known observed data, find the number of instances (or probability) satisfying the rule, find its average, and thereby evaluate the whole composite decision tree. Besides, it is also possible to estimate the stochastically most probable number of instances and distribution.
FIG. 27
1
2
Heretofore, an embodiment of the present invention has been described. The scope of the present invention is not restricted to the case where the explaining variables are the same with respect to object variables or decision trees. In other words, in the foregoing description, the case where the explaining variables are the same for respective object variables as shown in has been handled for brevity. However, the present invention can also be applied to the case where, for example, explaining variables for Y are different from explaining variables for Y.
If there are no duplications at all in explaining variables, however, the present invention can be applied, but the necessity of application is considered to be low. One of objects of the present invention is to implement inverse calculation for finding values of explaining variables that make a plurality of object variables desirable values. If explaining variables for object variables are completely different, there is no difference in processing contents irrespective of whether inverse calculations are performed independently using individual decision tree without combining the decision trees, or whether the decision trees are combined and then inverse calculation is performed. On the other hand, if there are partial duplications in explaining variables, the effect of the present embodiment is obtained.
Furthermore, in the present embodiment, an example in which two decision trees are combined has been described for brevity. Even if there are three or more decision trees, however, the present invention can be applied.
The above-described decision tree combination apparatus can be constructed by hardware. As a matter of course, however, the equivalent function can also be implemented by using a program.
Heretofore, the decision tree combination method and the decision tree improvement method have been described. Typically, the following advantages can be obtained by generation of the decision tree and data analysis using the decision tree.
FIG. 27
6
6
1
2
6
Generalization of the model and knowledge is facilitated by generating a decision tree from observed data. If a continuous value is used as a value of a variable, there is an advantage that moderate discrete is performed. In addition, since explaining variables that exert an influence upon the object variable, i.e., important explaining variables are automatically extracted when generating a decision tree, important explaining variables can be found. For example, in the data shown in , there is an explaining variable X. However, the explaining variable X is not present in the decision tree and the decision tree . Therefore, it can be the that the explaining variable X is not important. The decision tree is an effective model also in a sense that it provides the user with knowledge concerning data. Furthermore, the decision tree can cope with unknown data preferably while preventing excessive conformity to already known data.
According to the present embodiment, a plurality of decision trees are combined to generate a decision tree, which infers values of a plurality of object variables simultaneously on the basis of values of explaining variables, as heretofore described. By using this decision tree as an object decision tree in the first to fifth embodiments, therefore, inverse calculation for finding a condition to make a plurality of object variables simultaneously desirable values can be performed simple. If the combination method 1 is used as the decision tree combination method, then it suffices to add simple post-processing (a simple program) after generation of decision trees respectively for object variables, and consequently the processing is easy. In the combination method 2, a concise (easy to see) decision tree can be generated. In the combination method 3, a decision tree that is easy to understand correspondence to the original decisional tree can be generated, and the algorithm is also simple.
According to the present embodiment, a model with high precision can be constructed even if a loss value (a loss value of an object variable) is included in observed data. In the method of constructing a decision tree by regarding a direct product of object variables as one object variable (the second method described at the beginning of the present embodiment), there is a problem that if there is a loss value of an object variable in observed data the data of that portion cannot be used for construction of a decision tree and the precision of the constructed model falls. On the other hand, in the present embodiment, a decision tree for each object variable is first constructed. Thereafter, a composite decision tree is generated by combining the decision trees. In the present embodiment, therefore, a model (composite decision tree) with high precision can be constructed even if there is a loss value of an object variable in observed data. | |
QUESTION:
Given two strings A and B, find the minimum number of times A has to be repeated such that B is a substring of it. If no such solution, return -1.
For example, with A = “abcd” and B = “cdabcdab”.
Return 3, because by repeating A three times (“abcdabcdabcd”), B is a substring of it; and B is not a substring of A repeated two times (“abcdabcd”).
Note: The length of A and B will be between 1 and 10000.
EXPLANATION:
题目大意:求出最小的重复A数量,能够包含B,没有就返回-1. 考虑下,大概会分成两种情况, 1.既B是从A的第一个字符开始,这样只要长度大于等于B的长度就可以 2.B是从Ade中间字符开始,那么就需要再进行一次添加。
SOLUTION: | https://gaozhipeng.me/posts/686-Repeated-String-Match/ |
FIELD OF THE INVENTION
This invention relates generally to a thermocouple extension cable used in connecting a Type K or 20 Alloy/19 Alloy Thermocouple sensor to associated instrumentation and in particular, to a compensating extension cable comprising copper as the positive extension wire and a low nickel/high copper alloy as the negative extension wire which achieves the same accuracy limits as a standard Type K extension cable, but with significant material cost savings.
BACKGROUND OF THE INVENTION
An important parameter in many control systems is temperature. One of the most commonly employed mechanisms for dealing with the control of temperature is the thermocouple sensor. Thermocouple sensors are utilized to measure the temperature in high temperature environments such as those associated with autoclaves, furnaces, boilers, etc. Consequently, the prior art is replete with patents describing thermocouple devices of various configurations and constructions.
The Type K thermocouple sensor (Ni/10 Cr versus Ni/5 (Si, AI)) is presently employed in a wide array of temperature measurement and control applications. As stated earlier, the thermocouple sensor is coupled to the instrumentation by way of an extension cable. It is necessary that the thermal EMF of the extension cable is the same as the thermocouple sensor from 0° C. to the temperature of the transition where the extension cable is connected to the thermocouple sensor. It is desirable, from the standpoint of maintaining accuracy of measurement, for the thermocouple extension cable to exhibit the lowest possible loop resistance. Lowering the loop resistance of an extension cable allows the same instrument error limits with extended lengths of the extension cable. This is an advantage in applications where very long distances on the order of 100 feet or more exist between the thermocouple sensor and the instrumentation. For example, very long extension cables are employed between thermocouple sensors used in oil fields and the requisite instrumentation. These cables can be on the order of 100 feet or longer. Thus, in this application, a cable having lower loop resistance would greatly increase the accuracy of the temperature measurements.
Further, an extension wire that has a lower loop resistivity value allows the use of a smaller diameter wire for a given length of cable. Reducing the cable diameters also provides the benefit of enhanced cable flexibility.
Two standards setting forth the initial accuracy requirements for thermocouple sensor extension wire are maintained in the industry, one being the U.S. standard and the other being the international standard. The U.S. standard tolerance (established by ANSI, ISA, NIST, and others) for Type K extension wire (KX) is .+-.2.2° C. The IEC international standard tolerance for Type KX is .+-.2.5° C. In the U. S. standard, only type K thermocouple alloy is used as KX extension wire. The applicable temperature range for KX wire, both under the U.S. and the international standard is 0° to 200° C.
Most thermocouple extension cables are insulated with a low temperature material such as Poly Vinyl Chloride (PVC). The inventors herein have, therefore, recognized that PVC insulated KX cables provide an effective operating temperature well below 200° C. (The operating temperature of a PVC insulated KX cable is limited by the PVC insulation which has a maximum operating temperature 105° C. as established by Underwriters Laboratory [UL]). The consequences of all this is that the users are paying for unneeded accuracy above 105° C.
Hence, it becomes apparent from this disclosure that a switch to an extension cable manufactured from a metal costing less than KX, which meets the industry's initial accuracy requirements for thermocouple sensor extension cables up to 105° C., would result in a substantial cost savings to users of the cables. Moreover, if this metal also exhibits a lower loop resistivity and thus a lower loop resistance, this would allow the use of a smaller diameter wire to achieve additional material savings.
It is, therefore an object of the present invention to provide an alloy composition for use in the manufacture of thermocouple extension cables having a lower loop resistivity and lower material cost than presently available compositions used in the manufacture of KX extension cables for use up to 105° C.
SUMMARY OF THE INVENTION
An alloy composition used in the manufacture of the negative leg of an extension cable, comprises by weight 25.00% to 45.00% of nickel, 0. 10% to 1.75% of cobalt, 0.10% to 1.00% of manganese, less than 0.50% of iron, and the balance being of copper. A thermoelement, of a thermocouple extension cable, manufactured from this composition exhibits a resistivity of generally less than 300 ohms per circular rail foot. Hence, the loop resistivity of the cable, where the other thermoelement is made from copper (with a resistivity of 10 ohms/circular mil foot), is generally less than 310 ohms per circular mil foot and the calibration accuracy of the cable over the range of 0. degree. C. to 100° C. is within .+-.2.5° C. The preferred range of each of the elements in the alloy composition of the present invention includes 29.00 to 33. 00% of nickel, 0.30 to 1.00% of cobalt, 0. 10 to 0.70% of manganese, and less than 0.10% of iron.
In a preferred embodiment of the invention the alloy composition comprises by weight 30.00% of nickel, 69% of copper, 0.40% of manganese, 0.60% of cobalt. A thermoelement, of a thermocouple extension cable, manufactured from the preferred composition exhibits a resistivity of 240 ohms per circular mil foot. Hence, the loop resistivity of the cable, where the other thermoelement is made from copper, is 250 ohms per circular mil foot and the calibration accuracy of the cable 0° to 100° C. is within .+-.2.2° C.
BRIEF DESCRIPTION OF THE DRAWING
FIG. 1 illustrates a simple schematic circuit of an exemplary thermocouple arrangement employing thermocouple extension cables manufactured from the alloy composition of the present invention.
DETAILED DESCRIPTION OF THE INVENTION
Referring now to FIG. 1 there is shown a simple schematic circuit of an exemplary thermocouple arrangement employing thermocouple extension cables manufactured from the alloy composition of the present invention used as an extension for a Type K thermocouple sensor. Thermocouple 10 comprises a positive thermoelement 14 and a negative thermoelement 16. A sensing junction 12 is formed at the junction of thermoelements 14 and 16. The opposite ends of the thermoelements 14 and 16 form the intermediate junction 18 of thermocouple 10. Thermoelements 14 and 16 are coupled to the input of a high impedance operational amplifier 20 via a thermocouple extension cable 22 comprising a copper thermoelement 24 and a thermoelement 26 manufactured from the alloy of the present invention. The copper thermoelement 24 is coupled between the positive thermoelement 14 and the first input of amplifier 20 and the alloy thermoelement 26 is coupled between the negative thermoelement 16 and the second input of the amplifier 20. The output of the amplifier 20 is coupled to the input of a temperature detection circuit and display. The two inputs of Amp 20 constitute the reference junction of the thermocouple assembly shown.
It is understood that the arrangement illustrates one of many applications one of ordinary skill in the art would recognize for extension cables made from the alloy composition of the present invention. Further, the alloy composition of the present invention is intended for the manufacture of thermocouple extension cables, however, the alloy composition of the present invention is capable of being used in other applications where gains would result from the resistivity and material cost characteristics of the present invention.
As previously stated, copper versus the alloy of the present invention meets Type KX international specification of +2.5° C. from 0. degree. to 100° C. This cable composition presently enjoys a cost advantage when compared with KX as nickel is approximately four times more expensive than copper.
Furthermore, the inventors herein have recognized that the loop resistivity of an extension cable made from copper versus the present invention is generally half that of KX extension cable. As already mentioned, the accuracy of measurement with regard to the extension cable at the instrumentation is dependent on the loop resistance of the cable. As is well known, loop resistivity is defined as the combined resistivities of the positive and the negative thermoelements in ohms per circular mil foot. Further, the resistance of a conductor (in this case the thermoelements of the extension cable) can be expressed by the following equation:
R=p l/A
where R=resistance in ohms
l=length in feet
A=cross sectional area in square inches
p=resistivity in ohms per circular mil foot
For example, the resistivity of copper at ambient temperature (20. degree. C.)is 10 ohms per circular mil foot. Thus, a copper wire that is 1 foot in length and 0.001 inches in diameter will have a resistance at ambient temperature of 10 ohms. The resistivity of the conductor is dependent on its material characteristics, not its dimensions. The resistance can be obtained using the equation if p, land A are known.
Accordingly, it is proposed that a switch to copper versus the present invention from KX, would provide the same instrument error limit as KX with double the length of the extension cable because the loop resistivity of copper versus the present invention (310 ohms per circular mil foot) is generally half that of KX (600 ohms per circular mil foot). Consequently, using copper versus the present invention would allow a change to a smaller diameter wire to achieve material savings of approximately 50%.
In accordance with the present invention, the extension wire is formed from a metal alloy having a composition generally comprising copper (Cu), nickel (Ni), manganese (Mn), cobalt (Co), and a trace of iron (Fe). Since alloys are fundamentally an intentional mixture of two or more metals which are soluble in one another in the liquid state, alloying takes place by melting together the desired metals. As is well known, when the molten metals solidify, they remain soluble in one another or separate into intimate mechanical mixtures of the pure constituents metals.
Extension cables made from the alloy composition of the present invention are manufactured and processed according to methods that are well known in the art. Accordingly, an extension cable made from the alloy composition of the present invention can be made by induction melting the above-mentioned metals added in percentages which will be discussed in detail below, in a 800 pound or 3,000 pound furnace and pouring the melt into molds for 750 pound ingots. The ingots are hot rolled to rods, having a diameter of for example, 0.25 inches. The rods are then descaled and cleaned. After descaling and cleaning, the rods are drawn to various sizes, e.g. 16 gauge (0.051 inches in diameter), 20 gauge (0.032 inches in diameter) for solid conductor pairs or 34 gauge (0. 0063) for stranded conductor pairs. After wire drawing and cleaning, the thermocouple extension wire is annealed and coated with an appropriate insulation. The wires are spooled into proper sized reels which are checked for calibration. The compositional ranges of the abovementioned elements are depicted below in Table 1.
TABLE 1
______________________________________
ELEMENT RANGE PERCENTAGE BY WEIGHT
______________________________________
Nickel 25.00 to 45.00%
Copper Balance
Cobalt 0.10 to 1.75%
Manganese 0.10 to 1.00%
Iron less than 0.50%
______________________________________
A thermoelement, of a thermocouple extension cable, manufactured from this composition where the amount of nickel is close to 45% exhibits a resistivity of generally less than 300 ohms per circular mil foot. Hence, the loop resistivity of the cable, where the other thermoelement is made from copper, is generally less than 310 ohms per circular mil foot, (10+ 300), and the calibration error of the cable from 0° to 100° C. is within .+-.2.5° C.
The preferred range percentages of the elements making up the composition of the present invention are listed below in Table 2.
TABLE 2
______________________________________
PREFERRED RANGE
ELEMENT PERCENTAGE BY WEIGHT
______________________________________
Nickel 29.00 to 33.00%
Copper Balance
Cobalt 0.30 to 1.00%
Manganese 0.10 to 0.70%
Iron LESS THAN 0.10%
______________________________________
The loop resistivity of a cable manufactured from this composition, where the other thermoelement is made from copper, is essentially 250 ohms per circular mil foot, (10+240), and the calibration error of the cable from 0° to 100° C. is within .+-.2.2° C.
In a preferred embodiment of the invention the alloy composition comprises by weight 30.00% of nickel, 69% of copper, 0.40% of manganese, 0.60% of cobalt and essentially no iron. Extension cables were made from 4 duplicate trial melts, each of these melts being made according to the preferred composition. The extension cables drawn from these 4 melts displayed calibration errors at 100° C. of -0.7° C., -1.0. degree. C., -0.4° C. and -0.8° C. which are well within the +2.2° C. U.S. standard and .+-.2.5° C. international standard. A thermoelement, manufactured from the preferred composition exhibited a resistivity of 240 ohms per circular mil foot. The loop resistivity of the cable, where the other thermoelement is made from copper, was 250 ohms per circular mil foot.
The resistivities of copper, the positive and negative thermoelements of KX (designated KP and KN) and the alloy of the present invention (where the amount of nickel is approximately 45%) and the preferred alloy composition of the present invention are set forth below along with the loop resistivity for extension cables made from the individual thermoelements for comparison in Table 3.
TABLE 3
_________________________________________________________________________ _
RESISTIVITY OF
LOOP RESISTIVITY OF
THERMOELEMENT
EXTENSION CABLE
OHMS/CIR. MIL FT.
OHMS/CHT. MIL FT.
_________________________________________________________________________ _
COPPER 10
PREFERRED COMPOSITION
240 250
COMLPOSITION 300 310
(≈45% NICKEL)
KP 425
KN 177
KX 602
_________________________________________________________________________ _
Currently, the most popular cable sizes are 16 gauge and 20 gauge. From Table 3 it is clearly evident that the composition of the preferred embodiment allows a 20 gauge (0.032 inch diameter) cable made from the preferred composition to replace a 16 gauge (0.051 inch diameter) made from KX with the same instrument error because the loop resistivity of copper versus the alloy of the preferred composition is approximately half that of KX. Moreover, the loop resistivity of copper versus the preferred composition of the present invention is 20% lower than that of copper versus the composition of the present invention having approximately 45% nickel.
While the invention has been particularly shown and described with reference to specific exemplary embodiments of the alloy composition, other alloy compositions may become apparent to those skilled in the art that do not depart from the spirit and scope of the present invention. Hence, the present invention is deemed limited only by the appended claims and the reasonable interpretation thereof. | |
Why students can get the right answers without understanding the concepts
A study has shown that problems requiring interpretation of diagrams or graphs are useful for measuring students’ understanding of chemical kinetics. The research also revealed differences in the way students approach algorithmic and pictorial problems.
Chemical kinetics is a difficult topic for many students. They need a good understanding of the underlying concepts and a firm grasp of mathematics.
Students’ understanding of kinetics is mainly assessed through numerical problems. However, good performance can mask poor conceptual understanding. This is because numerical problems are algorithmic. This means that, generally, students need to manipulate a formula to compute a numerical answer using a plug-and-chug method, and they can solve the problems without fully understanding the underlying concepts. This contrasts with conceptual problems, which are often pictorial in nature and require the interpretation of a diagram or graph.
Kinetic misunderstanding
The research team, from Indonesia and the UK, examined student performance variations in a series of paired problems in chemical kinetics. 335 first-year undergraduate chemistry students in Indonesia and the UK participated.
One question in each problem pair was algorithmic, the other pictorial. Each pair of problems examined the same fundamental concept, but presented information differently and required different information processing and analytical steps. For example, one problem pair examined the impact of students’ understanding of reactant order upon reaction rate. The algorithmic problem involved typical handling of reaction rate and reactant concentration data, using given reaction orders, to determine the rate equation. These kinds of question are seen in many post-16 curriculums. The pictorial problem encoded this information in microscopic particle diagrams.
The results showed that the students’ performance in algorithmic problems was generally higher than in the pictorial problems. For example, students were able to select and use appropriate equations to get the correct answer in the algorithmic problem relating to the concept of half-life. However, the pictorial problem revealed a poor understanding of the concept of half-life. Therefore, pictorial questions may be a better diagnostic tool for dealing with misconceptions.
Also, students’ incorrect answers were different for the algorithmic problems compared to the pictorial ones, suggesting that students approached the question types differently.
Teaching tips
- You can access the test instrument that was used in this study for free. Try it in your classroom.
- For chemistry topics with a large mathematical component, good student performance on algorithmic problems does not necessarily mean they have understood the underlying concepts well.
- Algorithmic problems test students’ ability to select the correct mathematical formula and use it to get a numerical answer, while more conceptual problems test students’ ability to make links between the various representations used in chemistry (ie Johnstone’s triangle).
- Use pictorial problems as a diagnostic tool to address misconceptions you may not otherwise have been aware of. | https://edu.rsc.org/education-research/teaching-tips-for-chemical-kinetics/4014585.article |
21. The owner of Bun & Run Hamburgers wishes to compare the sales per day at two different locations. The mean number of hamburgers sold for 10 randomly selected days at Northside was 83.55 with a population standard deviation of 10.50. For a randomly selected 12 days at Southside, the mean number of hamburgers sold was 73.80 with a population standard deviation of 10.20. We wish to test whether there is a difference in the mean number of hamburgers sold at the two locations using a 5% significance level. What is the value of the test statistic in this case?
22. The owner of Bun & Run Hamburgers wishes to compare the sales per day at two different locations. The mean number of hamburgers sold for 10 randomly selected days at Northside was 83.55 with a population standard deviation of 12.45. For a randomly selected 12 days at Southside, the mean number of hamburgers sold was 73.80 with a population standard deviation of 14.25. We wish to test whether there is a difference in the mean number of hamburgers sold at the two locations using a 5% significance level. What is the correct conclusion for this hypothesis test?
23. The owner of Bun & Run Hamburgers wishes to compare the sales per day at two different locations. The mean number of hamburgers sold for 10 randomly selected days at Northside was 83.55 with a population standard deviation of 10.50. For a randomly selected 12 days at Southside, the mean number of hamburgers sold was 73.80 with a population standard deviation of 10.20. We wish to test whether there is a difference in the mean number of hamburgers sold at the two locations using a 5% significance level. What is the correct conclusion for this hypothesis test?
24. A researcher randomly sampled 30 graduates of an MBA program and recorded data concerning their starting salaries. The sample comprised of 18 women whose average starting salary is R48 000, and 12 men whose average starting salary is R55 000. It is known that the population standard deviations of starting salaries for women and men are R11 500 and R13 000 respectively. The researcher was attempting to show that female MBA graduates have significantly lower average starting salaries than male MBA graduates. What is the value of the test statistic in this case?
25. A researcher randomly sampled 30 graduates of an MBA program and recorded data concerning their starting salaries. The sample comprised of 18 women whose average starting salary is R45 000, and 12 men whose average starting salary is R55 000. It is known that the population standard deviations of starting salaries for women and men are R11 500 and R13 000 respectively. The researcher was attempting to show that female MBA graduates have significantly lower average starting salaries than male MBA graduates. What is the value of the test statistic in this case?
26. A researcher randomly sampled 30 graduates of an MBA program and recorded data concerning their starting salaries. The sample comprised of 18 women whose average starting salary is R48 000, and 12 men whose average starting salary is R52 000. It is known that the population standard deviations of starting salaries for women and men are R11 500 and R13 000 respectively. The researcher was attempting to show that female MBA graduates have significantly lower average starting salaries than male MBA graduates. What is the value of the test statistic in this case?
27. A researcher randomly sampled 30 graduates of an MBA program and recorded data concerning their starting salaries. The sample comprised of 18 women whose average starting salary is R48 000, and 12 men whose average starting salary is R55 000. It is known that the population standard deviations of starting salaries for women and men are R14 000 and R13 000 respectively. The researcher was attempting to show that female MBA graduates have significantly lower average starting salaries than male MBA graduates. What is the value of the test statistic in this case?
28. A researcher randomly sampled 30 graduates of an MBA program and recorded data concerning their starting salaries. The sample comprised of 18 women whose average starting salary is R48 000, and 12 men whose average starting salary is R55 000. It is known that the population standard deviations of starting salaries for women and men are R11 500 and R11 000 respectively. The researcher was attempting to show that female MBA graduates have significantly lower average starting salaries than male MBA graduates. What is the value of the test statistic in this case?
29. A researcher randomly sampled 30 graduates of an MBA program and recorded data concerning their starting salaries. The sample comprised of 18 women whose average starting salary is R48 000, and 12 men whose average starting salary is R55 000. It is known that the population standard deviations of starting salaries for women and men are R11 500 and R13 000 respectively. The researcher was attempting to show that female MBA graduates have significantly lower average starting salaries than male MBA graduates. What is the p-value of the test in this case?
30. A researcher randomly sampled 30 graduates of an MBA program and recorded data concerning their starting salaries. The sample comprised of 18 women whose average starting salary is R45 000, and 12 men whose average starting salary is R55 000. It is known that the population standard deviations of starting salaries for women and men are R11 500 and R13 000 respectively. The researcher was attempting to show that female MBA graduates have significantly lower average starting salaries than male MBA graduates. What is the p-value of the test in this case?
31. A researcher randomly sampled 30 graduates of an MBA program and recorded data concerning their starting salaries. The sample comprised of 18 women whose average starting salary is R48 000, and 12 men whose average starting salary is R52 000. It is known that the population standard deviations of starting salaries for women and men are R11 500 and R13 000 respectively. The researcher was attempting to show that female MBA graduates have significantly lower average starting salaries than male MBA graduates. What is the p-value of the test in this case?
32. A researcher randomly sampled 30 graduates of an MBA program and recorded data concerning their starting salaries. The sample comprised of 18 women whose average starting salary is R48 000, and 12 men whose average starting salary is R55 000. It is known that the population standard deviations of starting salaries for women and men are R14 000 and R13 000 respectively. The researcher was attempting to show that female MBA graduates have significantly lower average starting salaries than male MBA graduates. What is the p-value of the test in this case?
33. A researcher randomly sampled 30 graduates of an MBA program and recorded data concerning their starting salaries. The sample comprised of 18 women whose average starting salary is R48 000, and 12 men whose average starting salary is R55 000. It is known that the population standard deviations of starting salaries for women and men are R11 500 and R11 000 respectively. The researcher was attempting to show that female MBA graduates have significantly lower average starting salaries than male MBA graduates. What is the p-value of the test in this case?
34. It is known that the population variances of final exam marks for first year statistics students at a particular South African university are 45.3 for female students and 52.1 for male students. Samples of 27 female and 31 male first year statistics students from the university are selected and the sample mean exam marks are calculated. For females, the sample mean mark is 57.3% and for males the sample mean mark is 55.4%. If we wish to test whether females have, on average, higher exam marks than males, what would the test statistic value of the hypothesis test in this case be?
35. It is known that the population variances of final exam marks for first year statistics students at a particular South African university are 45.3 for female students and 52.1 for male students. Samples of 27 female and 31 male first year statistics students from the university are selected and the sample mean exam marks are calculated. For females, the sample mean mark is 58.3% and for males the sample mean mark is 55.4%. If we wish to test whether females have, on average, higher exam marks than males, what would the test statistic value of the hypothesis test in this case be?
36. It is known that the population variances of final exam marks for first year statistics students at a particular South African university are 45.3 for female students and 52.1 for male students. Samples of 27 female and 31 male first year statistics students from the university are selected and the sample mean exam marks are calculated. For females, the sample mean mark is 57.3% and for males the sample mean mark is 56.4%. If we wish to test whether females have, on average, higher exam marks than males, what would the test statistic value of the hypothesis test in this case be?
37. It is known that the population variances of final exam marks for first year statistics students at a particular South African university are 45.3 for female students and 52.1 for male students. Samples of 27 female and 31 male first year statistics students from the university are selected and the sample mean exam marks are calculated. For females, the sample mean mark is 57.3% and for males the sample mean mark is 59.4%. If we wish to test whether females have, on average, higher exam marks than males, what would the test statistic value of the hypothesis test in this case be?
38. It is known that the population variances of final exam marks for first year statistics students at a particular South African university are 45.3 for female students and 52.1 for male students. Samples of 27 female and 31 male first year statistics students from the university are selected and the sample mean exam marks are calculated. For females, the sample mean mark is 52.3% and for males the sample mean mark is 55.4%. If we wish to test whether females have, on average, higher exam marks than males, what would the test statistic value of the hypothesis test in this case be?
39. It is known that the population variances of final exam marks for first year statistics students at a particular South African university are 45.3 for female students and 52.1 for male students. Samples of 27 female and 31 male first year statistics students from the university are selected and the sample mean exam marks are calculated. For females, the sample mean mark is 57.3% and for males the sample mean mark is 55.4%. If we wish to test whether females have, on average, higher exam marks than males, what would the p-value of the hypothesis test in this case be?
40. It is known that the population variances of final exam marks for first year statistics students at a particular South African university are 45.3 for female students and 52.1 for male students. Samples of 27 female and 31 male first year statistics students from the university are selected and the sample mean exam marks are calculated. For females, the sample mean mark is 58.3% and for males the sample mean mark is 55.4%. If we wish to test whether females have, on average, higher exam marks than males, what would the p-value of the hypothesis test in this case be? | https://www.examrace.com/Sample-Objective-Questions/Statistics-Questions/Hypothesis-Testing-For-Two-Populations-Part-2.html |
Late night art sessions with only the sound of my thoughts and the continuous scribbling of pencil on paper is what drove my mind away from the worries of life in quarantine. With only a single lamp shining light onto my blank paper waiting to be filled with suppressed emotions, I had only one goal in my mind.
Having felt that there was no one else to listen to my current emotions, I used my own hand to communicate my thoughts with my blank, black paper. Starting with the base of the sketch, the white colored pencil outline is what began the sketching process of my desired artwork.
Once my sketch came out to my liking, I picked the different color of skin tones I would be using for my palette. Red, yellow, brown and purple were the pencils that would define the skin tone I wanted, colored pencil shavings of those colors would be messily sprawled all over my desk in the midst of the rush of creativity I was currently having.
The color purple can bring a meaning of magic or an escape from reality, so for this portrait, I decided to hide the most important features of a face that lets people know the depth of their emotions, which would be the eyes. I draw and color purple luminous flowers on top of the eyes and shade with different hues of purple.
The next most important feature I move on to is the mouth and the message I wish to convey through it. By drawing normal lips with a needle and thread piercing through them, the thread piercing the mouth shut all the way to the edge of the cheek, I choose to show the confined emotions that I haven’t been able to freely express as I wished to.
With the hair, I choose to draw it in its normal flowy state, but also a moth sits at the top of the hair, representing hope of seeing a gleam of light in the future. I draw the moth in its dullness of shades of browns to fit with the rest of the tones in the drawing.
As for the background of the portrait, I draw a glowing moon with a white gel pen, along with numerous small stars surrounding the moon and the portrait. A moon that symbolizes literal darkness and change one wants to see in their life is what I felt I should end my drawing journey with.
A portrait of myself with hints of emotions I can’t fully physically express is what I created, and felt satisfied with once I finished. Once again, expression through art helped me face my own self through these unpredictable times. | |
Q:
Ternary Operator Limits
Let's say we have following if statement:
int a = 1;
int b = 2;
if(a < b) {
System.out.println("A is less than B!");
}
else {
System.out.println("A is greater or equal to B!");
}
I have been wondering that if ternary operator replaces if statement when if statement consists from one line of code in each sub-block (if and else blocks), then why above example is not possible to write like this with ternary operator?
(a < b) ? System.out.println("A is less than B!") : System.out.println("A is greater or equal to B!");
A:
You can only use ? : for expressions, not statements. Try
System.out.println(a < b ? "A is less than B!" : "A is greater or equal to B!");
Note: this is also shorter/simpler.
A:
Because it doesn't replace an if statement.
The ternary operator only works on expressions and not statements, itself being an expression.
Because it is an expression, it is evaluated rather than executed, and it has to return a (non-void) value. The type of that value is inferred from the types of the two optional expressions specified, and the rules are fairly complex, with some unexpected gotchas.
(So as a rule I only use ?: in the simplest situations to keep the code easy to read.)
| |
Field
Background
Summary
Description of the Drawings
Description
Examples
The disclosure relates to the field of fabrication, and in particular, to fabrication robots.
Robots perform a variety of tasks upon parts in a fabrication environment. These tasks may include drilling, installing fasteners, welding, etc. When fabricating large parts, it is not uncommon for multiple robots to work collaboratively at the same time. When a group of robots perform work together on a part, it remains important that the robots do not collide with each other or the part. Collisions may damage the robots or the part, which may result in costly or time-consuming repairs.
In order to prevent collisions between robots working on the same part, all robots working on the part halt whenever one robot encounters a malfunction. This strategy successfully prevents collisions, but also increases overall downtime. That is, when robots are working collaboratively, if one breaks down then its collaborator robots will not be able to continue working. When a larger number of robots work collaboratively on a part, the amount of downtime at the group dramatically increases. This is because the likelihood of a single robot within the group encountering a malfunction increases as the number of robots in the group increases. There is also a desire to transition robots away from multifunction end effectors that go offline when one function of the robot encounters an error.
Therefore, it would be desirable to have a method and apparatus that take into account at least some of the issues discussed above, as well as other possible issues.
Embodiments described herein provide enhanced techniques for controlling robots in a manner that prevents collisions, while also allowing robots to continue working on a part after one robot has encountered a malfunction or otherwise become unable to conform its operations with a predefined schedule. The malfunctioning robot is removed in order to prevent collisions, and remaining robots continue to work on the part in accordance with their original schedule for performing work. When a functioning robot replaces the malfunctioning robot, the functioning robot is placed where the malfunctioning robot is currently scheduled to be. Because the functioning robot is placed in an already scheduled location that has been determined to be collision-free, it may continue work intended for the malfunctioning robot without issue.
One embodiment is a method for coordinating operations of robots performing work on a part. The method includes assigning a group of robots to a part, initiating work on the part via the group of robots, determining that a robot within the group is unable to continue performing work at a first location of the part, removing the robot from the group while other robots of the group continue performing the work, adding a functioning robot to the group at a second location that the robot is scheduled to occupy, and continuing work on the part via the group of robots.
A further embodiment is a non-transitory computer readable medium embodying programmed instructions which, when executed by a processor, are operable for performing a method for coordinating operations of robots performing work on a part The method includes assigning a group of robots to a part, initiating work on the part via the group of robots, determining that a robot within the group is unable to continue performing work at a first location of the part, removing the robot from the group while other robots of the group continue performing the work, adding a functioning robot to the group at a second location that the robot is scheduled to occupy, and continuing work on the part via the group of robots.
Another embodiment is an apparatus for coordinating operations of robots performing work on a part. The apparatus includes a controller that is configured to assign a group of robots to a part, and an interface that is configured to receive updates indicating work performed on the part via the group of robots. The controller is configured to determine that a robot from the group is unable to continue performing work at a first location of the part, direct removal of the robot from the group while other robots of the group continue performing work on the part, add a functioning robot to the group at a second location that the robot is scheduled to occupy, and direct the group to continue work on the part via the robots.
Yet another embodiment is a system for coordinating operations of robots performing work on a part. The system includes a group of robots, a part that the group of robots is scheduled to perform work upon, a controller that is configured to subdivide the part into regions, generate a schedule indicating where and when the group of robots will perform work on the part, and assign the group of robots to a part, and an interface that is configured to receive updates indicating work performed on the part via the group of robots. The controller is configured to determine that a robot from the group is not conforming with a schedule of work at a first location of the part, direct removal of the robot from the group while other robots of the group continue performing work on the part, add a functioning robot to the group at a second location that the robot is scheduled to occupy, and direct the group to continue work on the part via the robots.
A still further embodiment is a method for coordinating operations of robots performing work on a part. The method includes generating a schedule indicating where and when a group of robots will perform work on the part, the schedule including paths for robots that avoid collisions when timed with other robots, assigning the group of robots to the part, initiating work on the part via the group of robots according to the schedule, sampling a progress of the group of robots as the group of robots performs work on the part, and adjusting a speed of the robots based on the determined progress.
assigning a group of robots to a part (206);
initiating work on the part via the group of robots (208);
determining that a robot within the group is unable to continue performing work at a first location of the part (210);
removing the robot from the group while other robots of the group continue performing the work (212);
adding a functioning robot to the group at a second location that the robot is scheduled to occupy (214); and
continuing work on the part via the group of robots (216).
1. A method for coordinating operations of robots performing work on a part, the method comprising:
2. The method of clause 1 wherein:
working on the part via the group of robots comprises performing adjustments to a kinematic chain at of each of the robots.
subdividing each region into multiple sections; and
coordinating movement of the group of robots along the sections in a manner that prevents the robots from operating at the same time in sections that are directly adjacent.
3. The method of any of the preceding clauses 1 or 2 wherein:
the part is subdivided into regions, and the method further comprises:
determining volumes occupied by different robots during the work based on the schedule;
comparing volumes occupied by different robots over time to detect potential collisions; and
reporting any potential collisions that were detected.
4. The method of any of the preceding clauses 1-3 further comprising:
identifying a schedule for the group of robots;
5. The method of any of the preceding clauses 1-4 further comprising:
detecting potential collisions between robots in the group based on current positions, speeds, and tasks of robots; and
reporting any potential collisions that were detected.
subdividing the part into regions that are contiguous portions of the part and that do not overlap; and
assigning each robot in the group to a different region of the part.
6. The method of any of the preceding clauses 1-5 further comprising:
7. The method of clause 6 further comprising:
moving each robot in the group through its assigned region such that robots in two adjacent regions do not occupy a shared edge of the two adjacent regions at the same time.
repairing the robot to turn the robot into the functioning robot; and
continuing work on the part with the group of robots while repairing the robot.
8. The method of any of the preceding clauses 1-7 further comprising:
9. The method of any of the preceding clauses 1-8 wherein:
the work is selected from the group consisting of: installing fasteners, drilling, welding, and inspecting.
10. The method of any of the preceding clauses 1-9 wherein:
the robot is a malfunctioning robot.
11. A portion of an aircraft assembled according to the method of any of the preceding clauses 1-10.
The present apparatus, medium and method is also referred to in the following clauses which are not to be confused with the claims.
assigning a group of robots to a part (206);
initiating work on the part via the group of robots (208);
determining that a robot within the group is unable to continue performing work at a first location of the part (210);
removing the robot from the group while other robots of the group continue performing the work (212);
adding a functioning robot to the group at a second location that the robot is scheduled to occupy (214); and
continuing work on the part via the group of robots (216).
12. A non-transitory computer readable medium embodying programmed instructions which, when executed by a processor, are operable for performing a method for coordinating operations of robots performing work on a part, the method comprising:
13. The medium of clause 12 wherein:
working on the part via the group of robots comprises performing adjustments to a kinematic chain at of each of the robots.
subdividing each region into multiple sections; and
coordinating movement of the group of robots along the sections in a manner that prevents the robots from operating at the same time in sections that are directly adjacent.
14. The medium of any of the preceding clauses 12 or 13 wherein the method further comprises:
the part is subdivided into regions, and the method further comprises:
identifying a schedule for the group of robots;
determining volumes occupied by different robots during the work based on the schedule;
comparing volumes occupied by different robots over time to detect potential collisions; and
reporting any potential collisions that were detected.
15. The medium of any of the preceding clauses 12-14 wherein the method further comprises:
16. The medium of any of the preceding clauses 12-15 wherein the method further comprises:
detecting potential collisions between robots in the group based on current positions, speeds, and tasks of robots; and
reporting any potential collisions that were detected.
subdividing the part into regions that are contiguous portions of the part and that do not overlap; and
assigning each robot in the group to a different region of the part.
17. The medium of any of the preceding clauses 12-16 wherein:
18. The medium of clauseim 16 wherein the method further comprises:
moving each robot in the group through its assigned region such that robots in two adjacent regions do not occupy a shared edge of the two adjacent regions at the same time.
repairing the robot to turn the robot into the functioning robot; and
continuing work on the part with the group of robots while repairing the robot.
19. The medium of any of the preceding clauses 12-18 wherein the method further comprises:
20. The medium of any of the preceding clauses 12-19 wherein:
the work is selected from the group consisting of: installing fasteners, drilling, welding, and inspecting.
21. The medium of any of the preceding clauses 12-20 wherein:
the robot is a malfunctioning robot.
22. A portion of an aircraft assembled according to the method defined by the instructions stored on the computer readable medium of clause 12.
According to a further aspect of the present medium, there is provided:
a controller (144) that is configured to assign a group of robots (112, 114, 116, 118) to a part; and
an interface (142) that is configured to receive updates indicating work performed on the part via the group of robots;
the controller is configured to determine that a robot from the group is unable to continue performing work at a first location of the part, direct removal of the robot from the group while other robots of the group continue performing work on the part, add a functioning robot to the group at a second location that the robot is scheduled to occupy, and direct the group to continue work on the part via the robots.
23. An apparatus for coordinating operations of robots performing work on a part, the apparatus comprising:
24. Fabricating a portion of an aircraft using the apparatus of clause 23.
a group of robots (112, 114, 116, 118);
a part (120) that the group of robots is scheduled to perform work upon;
a controller (144) that is configured to subdivide the part into regions, generate a schedule indicating where and when the group of robots will perform work on the part, and assign the group of robots to a part; and
an interface (142) that is configured to receive updates indicating work performed on the part via the group of robots;
the controller is configured to determine that a robot from the group is not conforming with a schedule of work at a first location of the part, direct removal of the robot from the group while other robots of the group continue performing work on the part, add a functioning robot to the group at a second location that the robot is scheduled to occupy, and direct the group to continue work on the part via the robots.
25. A system for coordinating operations of robots performing work on a part, the system comprising:
26. Fabricating a portion of an aircraft using the apparatus of clause 25.
According to a further aspect of the present apparatus, there is provided:
generating a schedule indicating where and when a group of robots will perform work on the part, the schedule including paths for robots that avoid collisions when timed with other robots (1004);
initiating work on the part via the group of robots according to the schedule (1008);
sampling a progress of the group of robots as the group of robots performs work on the part (1010); and
assigning the group of robots to the part (1006);
adjusting a speed of the robots based on the determined progress (1012).
27. A method for coordinating operations of robots performing work on a part, the method comprising:
28. The method of clause 27 further comprising:
subdividing the part into regions, wherein assigning the group of robots to the part comprising assigning each robot in the group to a different region.
determining that a robot within the group is unable to continue performing work at a first location of the part;
adding a functioning robot to the group at a second location that the robot is scheduled to occupy; and
continuing work on the part via the group of robots.
removing the robot from the group while other robots of the group continue performing the work;
29. The method of any of the preceding clauses 26 or 27 further comprising:
30. A portion of an aircraft assembled according to the method of any of the preceding clauses 26-29.
According to a further aspect of the present method, there is provided:
Other illustrative embodiments (e.g., methods and computer-readable media relating to the foregoing embodiments) may be described below. The features, functions, and advantages that have been discussed can be achieved independently in various embodiments or may be combined in yet other embodiments further details of which can be seen with reference to the following description and drawings.
FIG. 1
illustrates a fabrication system for a part in an illustrative embodiment.
FIG. 2
is a flowchart illustrating a method for coordinating operations of robots in a fabrication system in an illustrative embodiment.
FIGS. 3-5
illustrate removal and replacement of a malfunctioning robot in a fabrication system in an illustrative embodiment.
FIG. 6
is a diagram illustrating coordinated movements between robots in different regions of a part in an illustrative embodiment.
FIG. 7
is a diagram illustrating regions of a part that have been divided into sections in an illustrative embodiment.
FIG. 8
is a block diagram of a schedule in an illustrative embodiment.
FIG. 9
is a diagram illustrating a collision avoidance model for robots in a fabrication system in an illustrative embodiment.
FIG. 10
is a flowchart illustrating a method for coordinating operations of robots in a fabrication system in an illustrative embodiment.
FIG. 11
is a block diagram of a fabrication system in an illustrative embodiment.
FIG. 12
is a flow diagram of aircraft production and service methodology in an illustrative embodiment.
FIG. 13
is a block diagram of an aircraft in an illustrative embodiment.
Some embodiments of the present disclosure are now described, by way of example only, and with reference to the accompanying drawings. The same reference number represents the same element or the same type of element on all drawings.
The figures and the following description provide specific illustrative embodiments of the disclosure. It will thus be appreciated that those skilled in the art will be able to devise various arrangements that, although not explicitly described or shown herein, embody the principles of the disclosure and are included within the scope of the disclosure. Furthermore, any examples described herein are intended to aid in understanding the principles of the disclosure, and are to be construed as being without limitation to such specifically recited examples and conditions. As a result, the disclosure is not limited to the specific embodiments or examples described below, but by the claims and their equivalents.
The systems described herein allow labor that would be performed by a robot with a multi-function end effector to be distributed across robots with single function end effectors. This division of labor to single function end effectors significantly increases the number of robots in use, resulting in a desire for synchronized control of the robots in order to avoid collisions. Using single-function end effectors provides a technical benefit because it overcomes problems related to multifunction end effectors breaking down when one function of a robot breaks down. The use of single-function end effectors also enables easier replacement of robots.
FIG. 1
illustrates a fabrication system 100 for a part 120 in an illustrative embodiment. Fabrication system 100 comprises any combination of systems, components, or devices that are operable to utilize robots in order to perform work on part 120 (e.g., a metal part, fiber reinforced composite part, etc.). Fabrication system 100 has been enhanced to coordinate these robots so that work continues uninterrupted on part 120, even when one or more of the robots have malfunctioned or otherwise become unable to continue performing work at the part 120.
In this embodiment, fabrication system 100 includes a group 110 of robots 112, 114, 116, and 118. These robots perform work upon regions 132, 134, 136, and 138, respectively, of part 120 in accordance with directions from robot coordination system 140. The work performed by robots 112-118 in regions 132-138 may comprise drilling, installing fasteners, welding, gluing, inspecting, or other operations. Robots perform actions by controlling their kinematic chains, such as a kinematic chain 117 of robot 116.
Robot coordination system 140 directs the operations of robots 112-118 in order to direct work at part 120 in a timely manner, and also to prevent robots 112-118 from colliding with each other or with part 120. In this embodiment, robot coordination system 140 includes interface 142, which provides instructions and receives updates from robots 112-118 indicating their progress. Interface 142 may comprise a wired communication interface such as an Ethernet interface or Universal Serial Bus (USB) interface, a wireless interface compliant with Wi-Fi or Bluetooth standards, etc.
Controller 144 reviews instructions stored in memory 146 in order to direct the operations of robots 112-118. In embodiments where robots 112-118 include their own internal controllers for repositioning and performing work at part 120, controller 144 may receive error messages or other notifications from robots 112-118. Controller 144 also generates schedules for operating robots 112-118 in tandem. Controller 144 checks and/or revises those schedules to prevent robots 112-118 from colliding during operation. In one embodiment, controller 144 additionally engages in real-time collision checking when directing work, based on updates from robots 112-118. Controller 144 may be implemented, for example, as custom circuitry, as a hardware processor executing programmed instructions, or some combination thereof. Memory 146 may be implemented as a solid state storage device, hard disk, etc.
Controller 144 has been enhanced to continue directing robots to perform work on part 120 even in circumstances where a robot has become unable to continue performing work. That is, if a robot encounters a malfunction or otherwise becomes unable to perform work in accordance with a schedule, controller 144 may continue to direct remaining robots to perform work on part 120. When a functioning robot is acquired, the functioning robot may initiate work at a location where the malfunctioning robot would currently be if it had continued operating normally. The robots then continue operating in accordance with their original schedule. Since the original schedule has already been checked to ensure that no collisions are present, the rest of the schedule may be run without any concern of collisions.
FIG. 2
Illustrative details of the operation of fabrication system 100 will be discussed with regard to . Assume, for this embodiment, that part 120 has been placed among robots 112-118 and awaits the performance of work in order for fabrication to be completed.
FIG. 2
FIG. 1
is a flowchart illustrating a method 200 for coordinating operations of robots in a fabrication system in an illustrative embodiment. The steps of method 200 are described with reference to fabrication system 100 of , but those skilled in the art will appreciate that method 200 may be performed in other systems. The steps of the flowcharts described herein are not all inclusive and may include other steps not shown. The steps described herein may also be performed in an alternative order.
With part 120 in place, controller 144 subdivides the part 120 into regions 132-138 in step 202, and generates a schedule indicating where and when robots 112-118 will perform work on part 120 in step 204. Regions 132-138 are defined by controller 144 as contiguous portions of part 120, and may comprise regions of the same size, regions expected to take the same amount of time for work to be completed by a robot, etc. In one embodiment, the schedule generated by controller 144 comprises one or more Numerical Control (NC) programs for robots 112-118. Controller 144 confirms that running the schedule will not result in any predicted potential collisions, and proceeds to step 202.
In step 206, controller 144 assigns a group 110 of robots 112-118 to part 120. In one embodiment, this includes assigning each of robots 112-118 to one of regions 132-138 of part 120. For example, controller 144 may provide a different NC program to each robot, based on the region in which that robot is located.
In step 208, controller 144 initiates work on the part 120 via the group 110 of robots. For example, controller 144 may initiate operations at each of robots 112-118 in accordance with a predefined schedule that synchronizes movements of the robots in order to assure collision avoidance. In some embodiments, controller 144 awaits further updates from the robots indicating their progress. In further embodiments, controller 144 selectively pauses, delays, or speeds up work at certain robots depending on their progress, in order to ensure that the robots continue operating in synchrony as dictated by the schedule. For example, if one robot has completed an operation faster than expected, controller 144 may briefly pause that robot in order to ensure that the robot continues to conform with the schedule, which is known to be collision free.
Progress continues on the part 120 as robots 112-118 continue to perform work. At some point in time, one of the robots encounters a condition that prevents it from being able to continue performing work on the part 120 in accordance with the schedule. For example, a robot may detect that it has broken or worn-down tooling that requires replacement, a robot may move to an unexpected position and halt, a robot may fail to perform work owing to a positioning error, the robot may encounter a runtime error, a robot may be in need of maintenance, etc. In any circumstance, this condition causes the robot to either halt work or to continue working at an undesirably reduced rate.
Because the robots 112-118 are each expected to perform their tasks on schedule along a predefined path that has already been crafted to prevent collisions, the malfunctioning robot presents a problem. If the robot can return to performing work in accordance with the schedule, then this problem may be corrected and work may continue. However, if the robot cannot return to performing work in accordance with the schedule, continued operation of the robot results in an unknown risk of collision with other robots. To this end, it is desirable to remove the robot before a collision occurs. After the robot has been removed, the remaining robots may continue in accordance with the schedule without risk of collision. The robot may then be replaced at a later time.
FIG. 3
In step 210, controller 144 determines that a robot is unable to continue performing work at a first location in one of the regions of the part. As used herein, a robot may be considered to be "malfunctioning" if it cannot continue to perform work in accordance with its schedule (i.e., if it cannot continue to conform with its schedule of work). This may be caused due to an error at an end effector of the robot, an error at an actuator that moves the robot, an error at a controller of the robot, etc. This determination may be based on an update received from the robot indicating a positioning or other error, may be based on the robot not reporting a confirmation to controller 144, due to a notification from a technician, etc. For example, in , robot 116 may encounter a malfunction at first location 320 after work 310 has been performed on part 120.
Upon detection of the robot that is unable to continue performing work in accordance with the schedule, the robot is removed from the group (e.g., physically and functionally) in step 212 while other robots of the group continue performing the work. For example, the robot may be physically removed from part 120. The robot may then be repaired to become a functioning robot, or may be replaced with a functioning robot.
FIG. 4
After a functioning robot has become available, in step 214 the functioning robot is added to the group at a second location in the region that the malfunctioning robot is scheduled to occupy. The second location is the location that the malfunctioning robot would currently occupy if the malfunctioning robot had been able to continue performing work on the part. For example, as shown in , a functioning robot 400 is placed at second location 410. Second location 410 is distinct from first location 320, and one or more instances of work 420 between first location 320 and second location 410 have not yet been completed. However, because the functioning robot 400 is placed at the location that is expected by the schedule, no collision checking or alteration of the schedule is required. That is, the functioning robot takes up collaboration spacing with the other robots where the malfunctioning robot would have been if it had not dropped out/been removed.
FIG. 5
In step 216, the group 110 of robots continue work on regions 132-136 of part 120. This may proceed until the schedule determined by controller 144 has been completed. For example, as shown in , each robot may continue performing work across the part 120. Method 200 may be repeated to perform additional work on the part 120, or even to perform different kinds of work on the part 120.
Method 200 provides a technical benefit over prior techniques, because it enables robots to continue operating according to a predetermined schedule when they are working on a part. This reduces downtime when operating on the part, and reduces the amount of processing resources needed for collision avoidance checking.
FIG. 6
The following FIGS. provide additional details of scheduling and collision avoidance in illustrative embodiments. is a diagram illustrating coordinated movements between robots in different regions of a part 650 in an illustrative embodiment. Robot 600 performs work within region 652 along path 610, followed by path 620, path 630, and path 640. Meanwhile, robot 602 performs work within region 654 along path 612, followed by path 622, path 632, and path 642. This means that robot 600 and robot 602 remain on opposite sides of their respective regions during work. It also means that robot 600 and robot 602 do not perform work at a shared edge 660 of region 652 and region 654 (which are adjacent) at the same time.
FIG. 7
is a diagram illustrating regions of a part 750 that have been divided into sections in an illustrative embodiment. In this embodiment, part 750 includes region 760 and region 770. Region 760 is worked upon by robot 710, while region 770 is worked upon by robot 720. Region 760 is subdivided into section 762 and section 764. Region 770 is subdivided into section 772 and section 774. Robot 710 performs work starting in section 762, and then section 764. Meanwhile, robot 720 performs work starting in section 772, and then in section 774. According to this scheduling technique, robots continue to work in sections that are not adjacent, which reduces a likelihood of collision during fabrication. Controller 144 may therefore coordinate movement of the robots along the sections in a manner that prevents the robots from operating at the same time in sections that are directly adjacent.
FIG. 8
is a block diagram of a schedule 800 in an illustrative embodiment. In this embodiment, schedule 800 is provided in the form of instructions for an NC program that is operated by controller 144. The instructions are for operations performed by each of multiple robots. Upon transmitting instructions for a set of operations (e.g., set 1, set 2, set 3) controller 144 pauses until the robots have confirmed completion of the operations. Upon receiving a confirmation from each robot, the controller 144 proceeds to provide a next set of instructions from the schedule 800. In further embodiments, schedules may be distributed across the robots for independent operation, may be performed without pausing, or may be implemented in any other suitable fashion. Controller 144 therefore independently determines what functions are to be performed by each robot, and where those functions will be performed. Controller 144 also determines a path for each robot that avoids collisions by performing the work at known timings. In further embodiments, Controller 144 periodically samples the progress of each of the robots to ensure that the schedule is being maintained.
FIG. 9
FIG. 9
is a diagram illustrating a collision avoidance model for robots in a fabrication system in an illustrative embodiment. According to , each robot is modeled as a volume, such as volume 910 or volume 920. The path 912 of volume 910 and path 922 of volume 920 across part 900 is also modeled. When checking a schedule for potential collisions, Controller 144 determines volumes occupied by the robots during the work based on the schedule. Controller 144 then compares volumes occupied by different robots over time to detect potential collisions, and reports any potential collisions that were detected. If any potential collisions are detected, controller 144 may also flag the schedule as unacceptable.
In further embodiments, controller 144 detects potential collisions between the robots based on current positions, speeds, and tasks of the robots, and reports any potential collisions that were detected. For example, one robot may drill the holes in a region while another robot inspects drilled holes in another region, while another robot installs pins into drilled and inspected holes in yet another region, all the while avoiding collisions due to scheduling and real-time control. This form of real-time collision checking may be implemented as a supplement to the schedule-based collision checking discussed above. The real-time collision checking provides a technical benefit by preventing collisions that would otherwise occur when a robot is not positioned where a schedule expects the robot to be.
In the following examples, additional processes, systems, and methods are described in the context of a fabrication system for a part.
FIG. 10
is a flowchart illustrating a method 1000 for coordinating operations of robots in a fabrication system in an illustrative embodiment. Method 1000 includes controller 144 determining tasks to be performed by each of multiple robots in step 1002. This may include determining what tasks are to be by robots on a part (e.g., drilling, inspecting, etc.) and where those tasks are to be accomplished. In step 1004, controller 144 generates a schedule indicating where and when a group of robots will perform work on the part. The schedule includes paths for robots that avoid collisions when timed with other robots. Thus, controller 144 confirms that when operating in accordance with the schedule, movements of robots within the group are coordinated to prevent collision. In step 1006, controller 144 assigns the group of robots to the part. For example, if the part is subdivided into regions, controller 144 may assign a different robot in the group to each of the regions. In step 1008, controller 144 initiates work on the part via the group of robots according to the schedule. For example, controller 144 may provide instructions to the robots in a timed manner in order to ensure compliance with the schedule. In step 1010, controller 144 samples a progress of the robots as the group of robots performs work on the part. For example, this may include receiving input from each robot indicating its current status and/or location. In step 1012, controller 144 adjusts a speed of one or more of the robots in the group based on the determined progress. For example, if some robots are proceeding more slowly than expected, controller 144 may either speed up these robots or slow down other robots in the group in order to ensure that the schedule is adhered to.
FIG. 11
is a block diagram of a fabrication system 1100 in an illustrative embodiment. Fabrication system 1100 operates on part 1110, which is divided into regions 1112 and sections 1114. In this embodiment, fabrication system 1100 includes robots 1120-1 through 1120-4, which include rigid bodies 1124-1 through 1124-4 that are repositioned by actuators 1122-1 through 1122-4 within kinematic chains 1128-1 through 1128-4. End effectors 1126-1 through 1126-4 perform work at robots 1120-1 through 1120-4, such as drilling or installing fasteners. Fabrication system 1100 also includes robot coordination system 1150. Robot coordination system 1150 includes an interface 1152 that is coupled for communication with robots 1120. Controller 1154 manages the operations of robots 1120 based on input from interface 1152, and accesses memory 1156 to store schedules and collision avoidance logic.
FIG. 12
FIG. 13
Referring more particularly to the drawings, embodiments of the disclosure may be described in the context of aircraft manufacturing and service in method 1200 as shown in and an aircraft 1202 as shown in . During pre-production, method 1200 may include specification and design 1204 of the aircraft 1202 and material procurement 1206. During production, component and subassembly manufacturing 1208 and system integration 1210 of the aircraft 1202 takes place. Thereafter, the aircraft 1202 may go through certification and delivery 1212 in order to be placed in service 1214. While in service by a customer, the aircraft 1202 is scheduled for routine work in maintenance and service 1216 (which may also include modification, reconfiguration, refurbishment, and so on). Apparatus and methods embodied herein may be employed during any one or more suitable stages of the production and service described in method 1200 (e.g., specification and design 1204, material procurement 1206, component and subassembly manufacturing 1208, system integration 1210, certification and delivery 1212, service 1214, maintenance and service 1216) and/or any suitable component of aircraft 1202 (e.g., airframe 1218, systems 1220, interior 1222, propulsion system 1224, electrical system 1226, hydraulic system 1228, environmental 1230).
Each of the processes of method 1200 may be performed or carried out by a system integrator, a third party, and/or an operator (e.g., a customer). For the purposes of this description, a system integrator may include without limitation any number of aircraft manufacturers and major-system subcontractors; a third party may include without limitation any number of vendors, subcontractors, and suppliers; and an operator may be an airline, leasing company, military entity, service organization, and so on.
FIG. 13
As shown in , the aircraft 1202 produced by method 1200 may include an airframe 1218 with a plurality of systems 1220 and an interior 1222. Examples of systems 1220 include one or more of a propulsion system 1224, an electrical system 1226, a hydraulic system 1228, and an environmental system 1230. Any number of other systems may be included. Although an aerospace example is shown, the principles of the invention may be applied to other industries, such as the automotive industry.
As already mentioned above, apparatus and methods embodied herein may be employed during any one or more of the stages of the production and service described in method 1200. For example, components or subassemblies corresponding to component and subassembly manufacturing 1208 may be fabricated or manufactured in a manner similar to components or subassemblies produced while the aircraft 1202 is in service. Also, one or more apparatus embodiments, method embodiments, or a combination thereof may be utilized during the subassembly manufacturing 1208 and system integration 1210, for example, by substantially expediting assembly of or reducing the cost of an aircraft 1202. Similarly, one or more of apparatus embodiments, method embodiments, or a combination thereof may be utilized while the aircraft 1202 is in service, for example and without limitation during the maintenance and service 1216. For example, the techniques and systems described herein may be used for material procurement 1206, component and subassembly manufacturing 1208, system integration 1210, service 1214, and/or maintenance and service 1216, and/or may be used for airframe 1218 and/or interior 1222. These techniques and systems may even be utilized for systems 1220, including, for example, propulsion system 1224, electrical system 1226, hydraulic 1228, and/or environmental system 1230.
In one embodiment, a part comprises a portion of airframe 1218, and is manufactured during component and subassembly manufacturing 1208. The part may then be assembled into an aircraft in system integration 1210, and then be utilized in service 1214 until wear renders the part unusable. Then, in maintenance and service 1216, the part may be discarded and replaced with a newly manufactured part. Inventive components and methods may be utilized throughout component and subassembly manufacturing 1208 in order to manufacture new parts.
Any of the various control elements (e.g., electrical or electronic components) shown in the figures or described herein may be implemented as hardware, a processor implementing software, a processor implementing firmware, or some combination of these. For example, an element may be implemented as dedicated hardware. Dedicated hardware elements may be referred to as "processors", "controllers", or some similar terminology. When provided by a processor, the functions may be provided by a single dedicated processor, by a single shared processor, or by a plurality of individual processors, some of which may be shared. Moreover, explicit use of the term "processor" or "controller" should not be construed to refer exclusively to hardware capable of executing software, and may implicitly include, without limitation, digital signal processor (DSP) hardware, a network processor, application specific integrated circuit (ASIC) or other circuitry, field programmable gate array (FPGA), read only memory (ROM) for storing software, random access memory (RAM), non-volatile storage, logic, or some other physical hardware component or module.
Also, a control element may be implemented as instructions executable by a processor or a computer to perform the functions of the element. Some examples of instructions are software, program code, and firmware. The instructions are operational when executed by the processor to direct the processor to perform the functions of the element. The instructions may be stored on storage devices that are readable by the processor. Some examples of the storage devices are digital or solid-state memories, magnetic storage media such as a magnetic disks and magnetic tapes, hard drives, or optically readable digital data storage media.
Although specific embodiments are described herein, the scope of the disclosure is not limited to those specific embodiments. The scope of the disclosure is defined by the following claims and any equivalents thereof. | |
The lightest stripes are actually white, but picture taken late at night...will try for an outdoor shot later and add it.
This is the 40th pair of socks I've knit, from the night I cast on the first one, January 28, 2012. I've learned a lot. I'm still learning. This pair has mostly Cascade 220 and Cascade 220 Superwash yarns in it, but the blue stripe at the top and the lighter blue stripe at the toe are both Bernat Sesame yarns from my mother's stash. The brighter yellow-gold (the Cascade 220 Superwash yarn) was purchased for socks for my son, but it turned out he wasn't that interested. The slightly darker caramel color stripes were on sale as a dyeing error and I thought the color looked "interesting." I still think it looks interesting and it will be used again. The yarn ends are still shaggy inside--not yet woven in--because I wanted a picture of them right after getting the toes together. Boosts my spirits.
Of the 40 pairs I've knitted, 3 were custom knit for friends, and one pair knit for myself was given to another friend for her husband to try on, and ended up going home with her daughter. A number have worn out. This is the 10th pair of shorty socks (all still whole) and there are 19 left in service of the 26 regular pairs knitted and kept for myself.
Pairs completed this year (including 2 pair started last year): 6 regular, 3 shorty. 4 pairs of regular socks (1 for friend) and all the shorty pairs were begun and finished so far this year. I need to keep up the pace, as some of the oldest socks are at or over their durability limit. | https://e-moon60.livejournal.com/483693.html |
Graduates Attending Higher Ed.
Class Size by Race/Ethnicity
Student Discipline
Student Discipline Days Missed
Arts Coursetaking
Grade Nine Course Passing
Advanced Course Completion
Digital Literacy and Computer Science Coursetaking
Related Links:
Attrition
2021-22 Attrition Report
This report provides the percentage of attrition by grade from the end of one school year to the beginning of the next for students enrolled in public schools, including charter schools, in the state. The information is as of October 1 of the school year selected.
More about the data.
Click any column header to sort ascending (first click) or descending (second click).
Student Group
K
1
2
3
4
5
6
7
8
9
10
11
All
All Students
7.7
14.3
20.0
12.5
6.7
89.5
28.3
Female
28.6
14.3
9.1
12.5
100.0
30.4
Male
0.0
0.0
25.0
0.0
81.8
26.1
High needs
14.3
16.7
16.7
0.0
0.0
77.8
23.4
Students w/ disabilities
20.0
Afr. Amer./Black
33.3
Hispanic/Latino
Multi-race, Non-Hisp./Lat.
29.4
White
10.0
12.5
10.0
0.0
0.0
92.3
24.6
Econ. Disadvantaged
12.5
0.0
85.7
25.7
A blank value indicates that either:
The school or district is new in the year selected
The school or district has no students enrolled in that grade level in the year selected
The school or district has no grade in the year selected for students from the previous year to advance
The data is suppressed because the enrollment total is less than 6.
A value of zero indicates that there was no attrition in that grade for the year and student group selected. | https://profiles.doe.mass.edu/attrition/default.aspx?orgcode=03000000&fycode=2022&orgtypecode=5&&dropDownOrgCode=2 |
To calculate annual averages, we analyzed data of 364 days (99.73% of year).
If in the average or annual total of some data is missing information of 10 or more days, this is not displayed.
The total rainfall value 0 (zero) may indicate that there has been no such measurement and / or the weather station does not broadcast.
|Data||Value||Computed days|
|Average annual temperature:||13.2°C||364|
|Annual average maximum temperature:||19.2°C||364|
|Average annual minimum temperature:||7.0°C||364|
|Annual average humidity:||70.7%||362|
|Rain or snow precipitation total annual:||639.85 mm||361|
|Annual average visibility:||-||-|
|Annual average wind speed:||7.7 km/h||364|
Number of days with extraordinary phenomena.
|Total days with rain:||109|
|Total days with snow:||7|
|Total days with thunderstorm:||29|
|Total days with fog:||44|
|Total days with tornado or funnel cloud:||0|
|Total days with hail:||3|
The highest temperature recorded was 36.8°C on 22 July.
The lowest temperature recorded was -8°C on 2 February.
The maximum wind speed recorded was 74.1 km/h on 1 April. | https://en.tutiempo.net/climate/2015/ws-160660.html |
The article discusses the principle of operation of a pass-through and a rocker switch, provides connection diagrams for switches designed to control lighting from two, three or more locations. Tips for the correct installation work related to the connection of the feed-through switches are given.
The idea behind the creation of a pass-through switch is not new, the first circuits appeared in the houses of radio amateurs back in the 60s, and it gained particular popularity in the 90s, when the first imported switches appeared on the market, “sharpened” under the control of a lamp from different places.
The device and principle of operation of the pass-through switch
The simplest representative of the family of pass-through switches is its one-key version.
Outwardly, it is no different from a conventional switch, except for the internal circuit, which is usually indicated on the back of the case..
The principle of operation of the pass-through switch is simple: when you move the switch button, the internal movable contact opens one circuit and automatically closes the second (the so-called changeover contact). In the figure, terminal “2” is a common contact, terminals “3” and “6” are a changeover.
The schematic diagram of the pass-through switch looks like this:
Using this effect, you can create the simplest circuit breaker circuit, in which one luminaire will be controlled from two different places at once:
1,2 – pass-through switches; 3 – to the lamp body
Connecting a pass-through switch
Installation is carried out with a three-core cable. To simplify installation work, its cores should be factory-color-coded. The section of the selected wire must withstand the load connected through it. Since the power of the switch contacts is limited to 10-16 A, copper flexible cable with a wire cross-section of 1 to 1.5 mm is most often used for installation.2.
How to connect a pass-through switch:
- On the pass-through switch, you need to find a common terminal (in the diagram it is indicated by the number “1”).
- On the first switch, located closest to the junction box, we bring the “phase” and connect it to the common terminal “1”. For installation, we use the brightest wire (usually red or orange, white is used in the explanatory picture).
- We land the two remaining wires on the output terminals of the pass-through switch (according to the scheme, these are terminals “2” and “3”), remember the correspondence of the color of the used core and the marking on the terminal block of the pass-through switch.
- On the second switch, we connect the cable in the same way as the first (we strictly observe the color coding of the wires and the corresponding switch terminals).
- In the junction box, we connect the brightest wire (in the explanatory figure it is white), which came from the second pass-through switch with the lamp phase.
- The other two wires, in accordance with the color coding, are connected with a wire of the same color that came from the first switch (for example, green with green, blue with blue, etc.), in the explanatory picture, green and red wires are connected.
- We immediately connect the neutral and ground wires in the junction box with a cable of the same purpose that goes to the lamp.
- We tighten the twists, if necessary, tinker, insulate the bare sections of the wires with high quality.
You can also use the following connection:
1 – branch box; 2 – to the luminaire body; 3, 4 – socket boxes
The assembly of the pass-through switch is performed in the following sequence:
1. We disassemble the switch.
2. We connect the wires to the pass-through switch, according to the diagram.
3. Insert the switch into the back box and fix it in it.
4. We close the switch with protective and decorative covers.
Important! Using the control, make sure which wire is “phase” in the junction box. Disconnect the power supply before carrying out installation work. Do not strand copper and aluminum wires together.
Checking the operation of the circuit
Ensure that each switch can both turn on and off the lamp, regardless of the position of the other switch..
Each switching of the pass-through switch should cause the electric lamps to turn off or on, if this does not happen, it is necessary to find and eliminate an error in the installation.
Two-key pass-through switches
These pass-through switches physically consist of two single pass-through switches assembled in one housing.
1 – two-key pass-through switch; 2 – pass-through switches
A double pass-through switch allows you to control multiple lamps at once. To do this, you need to collect the following scheme:
1, 2 – two-key pass-through switch; 3 – to the lamp body
For switching, you can use both three-core wires laid in parallel and six-core wires, the main thing is not to make a mistake when connecting.
The assembled circuit allows you to independently turn on and off two lamps or two lamps from two different locations.
For example, let’s turn on the lamp number 1 by changing the position of the first rocker switch.
Similarly, you can turn on the second lamp.
Disconnection can be performed using either the first or second switch.
Lighting control from three or more locations
In some cases, it is not enough to be able to control the lighting from two places. To effectively control the lighting of a three-story staircase, at least three control points are required. In this case, together with classic pass-through switches, an additional type of switch is used – a cross.
A cross switch is installed in the break of the connection between two pass-through switches, this allows you to create another lighting control point.
1, 3 – pass-through switches; 2 – cross switch; 4 – to the lamp body
Additional sequential installation of cross switches can increase the number of locations from which the lighting is controlled.
As you can see from the diagram, switching any of the switches will turn on or off the lighting..
Assembling the lamp control circuit from three different locations can be done as follows:
1 – pass-through switch; 2 – cross switch; 3, 5 – socket boxes for walk-through switches; 4 – socket for a cross switch; 6 – branch box; 7 – to the lamp body
Installation is carried out similarly to the above option with a single pass-through switch; installation requires a two and three-wire cable.
As can be seen from the material considered, with the help of pass-through switches, it is possible to organize the control of one lamp from two different places. The use of a cross switch allows the number of control points to be increased to three or more. | https://fin-radom.com/building/engineering-systems/how-to-connect-a-pass-through-switch.html |
Whodunit?
This can be a really fun mission or a really tedious mission or a really fast mission, depending upon how you'd like to handle it. There is a party at Summitmist Manor in Skingrad (it's just a little way in from the east gate, kind of across the street from your house). You're the (un)invited guest and everyone in the house is your target. How you take them out is entirely up to you. It should preferably be done quietly and privately (which won't be hard to do), but it really doesn't matter.
The guests are at Summitmist Manor to find a chest of gold hidden in the house (there isn't one, so don't bother looking). They are a strange assortment right out of an Agatha Christie novel:
If you want to do it really fast, talk to the doorman (who will give you the key to let you out again), go in, kill each of the guests whenever you can get them alone and leave. Aside from the few golds each guest has (and Neville's armor), there isn't much in the house worth looting.
The reason why it's both fun and tedious is that the way each guest reponds to the situation and other guests depends upon how many people are currently alive and their Disposition toward you, so there are... let's see... five guests times three different Disposition levels (less than 30, 30-70, and more than 70), times five times four times three times two equals ... ummm... a lot of different ways that it can play out (break out a calculator if you're really interested). I am reminded of a "Babylon 5" episode in which an assassin is sent to take out someone with instructions that are someting to the effect of "you are to know pain and then you are to know fear and then you are to die." This is about like that. If you play your cards right, you can get them to kill each other for you. You'll probably have to take out the last guest on your own since it may be kind of obvious who the killer is by that point. But if you played your cards right, the last surviving guest will think that they have killed the killer and are now safe.
You can also play "fun with Poisoned Apples" if you'd like. Get some of those yummy Poisoned Apples from M'raaj-Dar. When you get into the house, collect all of the food that you can find and put the Poisoned Apples out on the table. Then have a seat and wait for people to get hungry. It might take a day or two, but everyone dies and you might even get the last ones to kill each other so you'll have leftovers for later.
Report back to Ocheeva when the last guest is dead. Your reward is your usual fee and, if you weren't discovered as the assassin, the "Night Mother's Blessing" (permanent boosts to Sneak, Blade, Security, Acrobatics and Marksman skills). | https://tesguides.com/tes4/factions/darkbrotherhood/whodunit.htm |
In bacteria and archaea, CRISPR-Cas constitutes an adaptive RNA-mediated defence system which targets invading phages or plasmids in three steps: adaptation via integration of viral or plasmid DNA-derived spacers into the CRISPR locus; expression of short guide CRISPR RNAs (crRNAs) consisting of unique single repeat-spacer units and interference with the invading cognate foreign genomes. Short mature crRNAs are key elements in the interference step of the immune pathway. CRISPR-Cas systems are categorized into three major types (I, II and III) and several subtypes. Types I (except for I-C) and III employ the Cas6-mediated crRNAs processing while the type II CRISPR-Cas systems use the tracrRNA-guided processing mechanism with endogenous RNase III. Type I-C is associated with the Cas5d-mediated processing.
crRNA maturation is an important process of CRISPR immunity. CRISPR arrays are comprised of a set of cas genes and identical repeats interspaced by spacer sequences acquired from invading mobile genomes. The repeat-spacer array is transcribed as a precursor crRNA (pre-crRNA) molecule. This pre-crRNA undergoes one or two maturation steps to generate the mature crRNAs that function as guides in destruction of invading DNA or RNA. Generally, primary cleavage of the pre-crRNA occurs at a specific site within the repeats to yield crRNAs that consist of the entire spacer sequence flanked by partial repeat sequences. In some cases, an additional secondary cleavage step is required to generate the active mature crRNAs. The maturation of the crRNAs is critical for the activity of the system.
Fig 1. Pathways and enzymes of CRISPR RNA processing. A. Type I and III CRISPR-Cas systems; B. Type II CRISPR-Cas system
A common theme among the CRISPR-Cas types is the transcription of the pre-crRNA and the first muturation processing event within the repeats. In Types I and III, a protein of the Cas6 family or alternatively Cas5d catalyzes this step. The processed crRNAs from Types I-C, I-E and I-F do not undergo further maturation, whereas in at least Types I-A, I-B and I-D, as well as in Types II and III, a second maturation step produces the active crRNAs. In Types I and III, cleavage within each repeat by Cas6 releases the spacers bearing portions of the repeat on its 5' and 3' ends. The 5' flanking repeat of the crRNAs is the last 8 nts of the preceding repeat and the 3' flanking repeat is the remaining sequence of the downstream repeat. In some systems, the 3' flanking repeat sequences are further processed by uncharacterized exonucleases.
Each member of the Cas6 family of endoribonucleases recognizes a unique RNA sequence and collectively, Cas6 nucleases process a wide range of substrates of different sequences and secondary structures. Cas5d is the second distinct class of endoribonucleases responsible for processing crRNA in CRISPR system that lacks Cas6, such as Type I-C. Similar to Cas6, Cas5d also recognizes specific features of and cleaves within the repeat, resulting in crRNAs containing spacer sequences flanked by repeat sequences. While Cas5d has been shown functionally similar to Cas6, other subtypes of Cas5 have no roles in crRNA maturation. Rather, they play a key role in assembly of surveillance or effector complexes.
Two Cas proteins, Cas6 and Cas5d, have been identified as endoribonucleases that cleave within the repeat sequences of pre-crRNA to generate the mature crRNAs. However, their homologues are missing in many CRISPR-Cas subtypes, suggesting the existence of alternate crRNA maturation pathways involving other Cas proteins and/or fundamentally different RNA processing events. Another pathway of CRISPR activation has been uncovered. In Type II CRISPR-Cas systems, a unique crRNA maturation pathway which distinct from the Type I and III, involves the coordinated action of three novel factors: the trans-acting small RNA (tracrRNA), the host-encoded endoribonuclease RNase III and the Cas9 protein.
TracrRNA and pre-crRNA undergo coprocessing through the double-stranded substrate formed by the base pairing of tracrRNA anti-repeat and pre-crRNA repeats. And then, the endogenous RNase III-a general RNA processing factor in bacteria-is recruited to cleave the duplex RNA which was stabilized by the Cas9 protein, to generate predictable dsDNA breaks into the target sequence. RNase III cleaves both strands of the dsRNA with a two base pair separation, resulting in a cleavage intermediate further processed by the Cas9 class protein. RNase III which serves as a host factor in tracrRNA-mediated crRNA maturation, is an evolutionarily conserved endoribonuclease involved in many biological processes.
In the interference step, crRNAs combine with Cas proteins to form an effector complex which recognizes the target sequence in the invasive nucleic acid by base pairing to the complementary strand and induces sequence-specific cleavage, thereby preventing proliferation and propagation of foreign genetic elements. The structural organization and function of effector ribonucleoprotein (RNP) complexes involved in crRNA-mediated silencing of foreign nucleic acids differ between distinct CRISPR-Cas types. In the type I-E systems of E. coli and Streptococcus thermophilus CRISPR4-Cas, crRNAs are incorporated into a multisubunit effector complex called Cascade (CRISPR-associated complex for antiviral defense), which binds to the target DNA and triggers degradation by the signature Cas3 protein.
In type III CRISPR-Cas systems of Sulfolobus solfataricus and Pyrococcus furiosus, RNP complexes of Cas RAMP (Cmr) proteins and crRNA recognize and cleave synthetic RNA in vitro, while the CRISPR-Cas system of Staphylococcus epidermidis targets DNA in vivo. In type II CRISPR-Cas systems, as exemplified by S. pyogenes and S. thermophilus CRISPR3-Cas (St-CRISPR3-Cas), a single Cas9 protein, instead of a multisubunit Cascade or Csm/Cmr protein complex, provides DNA silencing. In fact, type II effector complex functions as an RNA-guided endonuclease which achieves target sequence recognition by crRNA and employs Cas9 protein for DNA cleavage within the target protospacer.
1. Deltcheva et al. CRISPR RNA maturation by trans-encoded small RNA and host factor RNase III. Nature. 2011 March 31; 471(7340): 602–607.
2. Karvelis et al. crRNA and tracrRNA guide Cas9-mediated DNA interference in Streptococcus thermophilus. RNA Biology. May 2013; 10(5): 841–851.
3. Haurwitz et al. Sequence- and structure-specific RNA processing by a CRISPR endonuclease. Science. 2010 September 10; 329(5997): 1355–1358.
4. Hong Li. Structural principles of CRISPR RNA processing. Structure. 2015 January 6; 23(1): 13–20.
5. Charpentier et al. Biogenesis pathways of RNA guides in archaeal and bacterial CRISPR-Cas adaptive immunity. FEMS Microbiology Reviews. fuv023, 39, 2015, 428–441. | http://tw.sinobiological.com/crispr-rna-crrna.html |
Poetry is a type of artistic expression in its own right. Did you realize, though, that there are over 50 different genres of poetry? Outside of advanced poetry workshops or in-depth studies, we prefer to concentrate on seven forms of poetry. Haiku, free verse, sonnets, and acrostic poems are all popular styles of poetry.
The best way to learn about poetry is by reading as much as possible. Try writing your own poems: you may find some new ideas how to use language more effectively.
In conclusion, poetry is the art of expressing thoughts and feelings in words. There are many different styles of poetry, but they all share one common theme - the imagination of the poet. No matter what kind of poetry you read, understand that it is meant to inspire and convey emotions.
It is hard to include every type of poetry used by writers, but the most popular are sonnets, Shakespearian sonnets, haiku, limericks, odes, epics, and acrostics. There are also villanelles, sestinas, pantunos, and tercets.
Poetry is defined as "the art of creating poems." That is not an exact definition, but it will serve for now. It is difficult to define poetry because there are so many different types of poems. Poems can be described as a series of words, phrases, or sentences that make up a work of literature.
In general, poems are divided into five-line sections called stanzas. Some poems have additional structural elements such as four-line verses, six-line quatrains, eight-line octaves, etc. A poem may have more than one stanza or any other special structure element.
The term "poem" usually refers to something that uses poetic language and forms. This includes lyrics, sonnets, villanelles, and so on. A poem can also be considered anything written in verse: novels, short stories, essays, and so forth. Finally, a poem can be interpreted by someone who reads it aloud; this includes songs, chants, and recitations.
Ballads, epics, idylls, and laments are the four basic forms of narrative poetry. Learn more about the three kinds of poetry.
Poetry is still thought to be the exclusive owner of the three major poetic forms: lyric, narrative, and dramatic. Each form can then be further subdivided into several subcategories, each with its own rhyme scheme, rhythm, and/or style. A song-like quality can be found in an expressive piece of literature focused on thinking and emotion. A story can be told through a series of events with a beginning, middle, and end. Last, a drama requires conflict between characters who engage in verbal sparring or role playing to express their views about something important to them.
These are only some examples of how poetry can be divided up. The main thing to remember is that there are many different ways of dividing up poetry, and it's possible to divide it into as many or few categories as you want!
For example, one common way of classifying poems is by subject matter. This could be done either by category or by sequence. A sequence is a group or list of poems written by the same author or artists. These can be autographs (i.e., poems written by one person) or epigrams (i.e., short poems written by multiple people). Categories are general groups of subjects covered in poems. Some examples include love poems, political poems, and religious poems. Category definitions can vary depending on who is assigning them to poems.
Poetry is further split into three genres: lyric, epic, and dramatic. All of the shorter forms of poetry, such as song, ode, ballad, elegy, and sonnet, are included in the lyrics. Comedy, tragedy, melodrama, and hybrids such as tragicomedy can all be found in dramatic poetry. Long narrative poems that use regular lines of iambic pentameter to form a continuous poem are called epics.
The term "form" can also be used to describe the general shape or structure of a work of literature. For example, many novels follow a basic plot outline that typically includes a beginning, middle, and end. Some novels, however, do not adhere to this formula and may instead focus on a particular subject over several episodes or characters who are not necessarily linked together linearly through time.
Literary forms can also include styles of writing. For example, the epic form is most commonly associated with ancient Greek and Roman poets but has been adapted by many other writers since then. Lyric poetry is known for its use of iambic verse, which consists of two pairs of metrically identical stressed syllables followed by an unstressed syllable. This pattern is usually repeated throughout the poem without any variation or exception.
Dramatic poetry is defined by the presence of drama or performance elements within the text itself. | https://authorscast.com/how-many-types-of-poetry-are-there |
horseshoes(redirected from Horseshoe pit)
Also found in: Thesaurus.
horse·shoe(hôrs′sho͞o′, hôr′sho͞o′)
n.
1. A flat U-shaped metal plate fitted and nailed to the bottom of a horse's hoof for protection.
2. A U-shaped object similar to a horseshoe.
3. horseshoes(used with a sing. verb) A game in which players toss horseshoes or horseshoe-shaped pieces at a stake so as to encircle it or come closer to it than the other players.
tr.v. horse·shoed, horse·shoe·ing, horse·shoes
To fit with horseshoes.
American Heritage® Dictionary of the English Language, Fifth Edition. Copyright © 2016 by Houghton Mifflin Harcourt Publishing Company. Published by Houghton Mifflin Harcourt Publishing Company. All rights reserved.
horseshoes(ˈhɔːsˌʃuːz)
n
(Games, other than specified) (functioning as singular) a game in which the players try to throw horseshoes so that they encircle a stake in the ground some distance away
Collins English Dictionary – Complete and Unabridged, 12th Edition 2014 © HarperCollins Publishers 1991, 1994, 1998, 2000, 2003, 2006, 2007, 2009, 2011, 2014
ThesaurusAntonymsRelated WordsSynonymsLegend:
Switch to new thesaurus
|Noun||1.||horseshoes - a game in which iron rings (or open iron rings) are thrown at a stake in the ground in the hope of encircling it|
leaner - (horseshoes) the throw of a horseshoe so as to lean against (but not encircle) the stake
ringer - (horseshoes) the successful throw of a horseshoe or quoit so as to encircle a stake or peg
outdoor game - an athletic game that is played outdoors
Based on WordNet 3.0, Farlex clipart collection. © 2003-2012 Princeton University, Farlex Inc. | https://www.thefreedictionary.com/Horseshoe+pit |
Stacked Cabinet No.5 is a piece that marks an important transition in the life of Dust furniture. In late 2007, Dust acquired a new piece of equipment, a CNC router, that allowed even more precision in the cutting and assembling of Vincent's designs. Stacked Cabinet No.5 integrates two pieces that are not only stacked, but also intersect & overlap. This intersection was made possible by the precision of the CNC router. You will see that pieces made after this point begin to display this freedom of creation in increasingly interesting ways (the functioning top drawer of Together We Can, the intersection of Stacked Cabinet No.7, the piece construction of the Mantel Clock, etc).
Choose from one of the shown colors or email us your color choice. | https://www.zinhome.com/stacked-cabinet-no-5/ |
# Wisconsin Highway 55
State Trunk Highway 55 (often called Highway 55, STH-55, or WIS 55) is a state highway in Wisconsin, United States. It travels south-to-north in the northeastern part of Wisconsin from an intersection with U.S. Route 151 (US 151) approximately 1.5 miles (2.4 km) north of Brothertown, near the eastern shore of Lake Winnebago in Calumet County, to the Michigan state line at the Brule River approximately one mile (1.6 km) northeast of Nelma in Forest County, where it connects to M-73.
## Route description
Along its route, STH-55 serves Kaukauna, Shawano, the Menominee Indian Reservation, Crandon, and the Nicolet side of the Chequamegon-Nicolet National Forest.
## History
A new roundabout was opened at the intersection of WIS 55 and US 10 between Sherwood and Kaukauna in the autumn of 2009. Another roundabout recently opened at the busy intersection of WIS 55 and WIS 114 approximately one mile (1.6 km) west of Sherwood, Wisconsin.
Over the summer of 2018, a 1 mile (1.6 km) section of the Wisconsin Highway 55 at its interchange with I-41 in Kaukauna was reconstructed and 4 new roundabouts were be added. | https://en.wikipedia.org/wiki/Wisconsin_Highway_55 |
So first of all, I just want to say a HUGE thank you to everyone who commented on my last post allowing me to write this one! You guys are awesome! This post is going to be packed with information on fairies, so I hope you're ready! ;)
First of all, I want to talk about the term fairy. What is a fairy? Something I recently learned is that it is much like saying animal. It's a very basic general term that leads to hundreds, even thousands of different species. Typically when I heard the term fairy I pictured a female humanlike figure with wings. The term "fairy" also spelled "faerie" is used to depict the denizens of the "otherworld" or other "realm". What I now know is that a leprechaun, a goblin, a troll, a satyr, an elf, a mermaid, a pillywiggin, a pixie, and a sprite are all fairies, to name only a few.
In my book Beyond The Veil, I mention several different kinds of fairies, including dryads, hobgoblins, and of course pillywiggins, pixies, and sprites. Below I have recorded a brief cyclopedia of the fairies to the best of my knowledge.
Pillywiggins.
BOOK EXCERPT:
Type: Pillywiggin
Name: Shaylee
Habitat: the fields
Pillywiggins live in the fields of wildflowers, for they are flower fairies. They tend to the flowers in people's gardens and in the fields.
KNOWLEDGE: Pillywiggins originated in England and Wales. These flower fairies are seen in spring and summer months and hibernate during the colder months. Pillywiggins typically live in fields of wildflowers, but tend to live in gardens as well. They are tiny, playful, loving, and excellent gardeners. They do not plant gardens, but rather, they tend to those gardens belonging to people. Though they are generally uninterested in humans, they sometimes show themselves to those who have a good heart and have asked to see them. After a human has asked to see a fairy, a pillywiggin will sometimes watch that person for several days to make sure that that person has a good heart. One thing to note is that pillywiggins are known to visit more often if there are lemon treats left out for them. Often pillywiggins are seen riding dragonflies or bees, as they have no wings of their own. Some believe that they have mated with insects which would explain their long antennas and green skin. Their clothes are made of flowers, and often they are seen wearing hats made of leaves or flower petals. Their weddings and dances and other parties (which they picked up from humans) are the most spectacular events to behold, as noted from Beyond The Veil:
"I know how taken you are with the village and the humans. A pillywiggin long before our time, a queen whose name is long forgotten, was not so different than yourself. She became so infatuated with their world and wished to see them up close. And so she left the fairy realm behind her. She saw a great many things. But more than anything, more than the people she had encountered or the vastness of the earth she saw or the sense of freedom she felt, she was stricken with awe by their celebrations and dances. It was like nothing our world had ever known. But she made it known. She returned home, took a husband for herself, and held a grand celebration for all sorts of fairy folk. It was the first wedding our world had known, with food and music and dancing."
Pixies.
BOOK EXCERPT:
Type: Pixie
Name: JuJuBee
Habitat: the Forest
Pixies are tiny human-like figures with large pointed ears, wild eyes filled with childish glee, and feet that carry them as fast as that of a hunted fox. Pixies are mischievous beings who delight in deliberately tormenting and misleading people, sometimes with the use of magic, sometimes not.
KNOWLEDGE: Pixies are thought to be of Celtic origin. They are very small and extremely fast. Contrary to common belief, pixies do not have wings, but they do have very pointed ears. They are very mischievous and delight in misleading and/or tormenting people; causing one to lose their way or walk in circles without realizing it, or even throwing one's clothes off the line. To avoid pixie magic, one must wear his or her clothes inside-out. Pixies typically wear clothes made of grass. It is believed that pixies are so playful because they are the souls of humans who died as infants. Fairy rings are the result of pixies dancing so it is said, and those who enter a fairy ring are transported into the fairy realm, instantly becoming invisible to the human world.
Sprites.
BOOK EXCERPT:
Type: Sprite
Name: Hollyhocks
Habitat: the Forest
Sprites are perhaps the most beautiful of all the fairies for reasons not only in appearance, but also in manner. They are very kind-hearted, caring, and innocent by nature. They resemble humans, except for their pale and almost translucent soft skin, and their elegant gossamer wings protruding from their backs. The fairies all lived at peace with one another for the most part, but unlike the hobgoblins and the pixies, the sprites avoided all human contact. The magic they worked was very important and like nothing else to behold. Nothing short than the work of a skillful artist, they painted the leaves in the autumn, carefully but swiftly, though not too swift; they enjoyed to make the autumn last before the winter came upon them and carried the leaves away in the wind.
KNOWLEDGE: Sprite comes from the Latin word "spiritus", meaning spirit or ghost. Originating from Celtic folklore, sprites were generally thought of as tiny ghost-like beings because of the faint glowing light coming from their bodies. Though they usually avoid human contact, they are very peaceful. Sprites are some of the most creative fairies and are well practiced in the arts of music, poetry, and painting. They were given one very important task, which is to change the color of the leaves in the fall. Sprites are known to vary in size, some say they are even capable of shape-shifting. These winged fairies are very kind-hearted and caring. Because of their innocent nature, they see it fit to wear no clothing as noted in Beyond The Veil:
You see Dear Reader, the sprites live in a different world than the humans entirely. They live in peace with one another always. They know no hate, only love. They are innocent in the ways of man and they are not ashamed to be naked. | http://www.vanessapaigeisrael.com/2015/11/ |
The alphabet song is something most children learn, but research shows that to become successful readers children also need to apprehend the “alphabetic principle”: the concept that letters are symbols that represent the sounds of speech. Children need to understand that letters not only make predictable sounds, but also correspond to the sounds of spoken language.
Here are some ideas to reinforce this concept:
- Prop: Cell Phone. Explain that if you want to communicate to another person you can do it a few different ways. You can call another person and speak into the phone. The person you are calling will hear your voice and understand your message. Or you can send the same message via email or text message. When you email or text, the letters you type into your cell phone represent the sounds of the words used in your voice message. If the person on the other end knows how to read your message, they’ll understand just from looking at the symbols (letters) in your message what you are communicating.
- Prop: A Book, Magazine or Newspaper. (Be sure to choose something with few or no pictures.) Explain to your child that you can read the “secret code” represented by all those letters on the page. The letters, words, and sentences represent the sounds of spoken language. Explain that while your child naturally and easily learned to speak, learning the “secret code” requires more effort. To learn to read, one must learn the name, shape, and sound(s) represented by each letter of the alphabet.
- Props: Sticky note, pen, and a picture book. Before you begin reading the picture book to your child, point to a single object pictured in the story and ask what it is. Be sure to choose an object that is easy to spell. (For example: cat, dog, dad, pig or another three letter word.) Write that word on the sticky note and explain that you are using letters to represent the word your child just said. Place the sticky note in the book on the page where the object appeared. Then, read the book to your child. When you come to the sticky note ask your child to read it to you.
Learning to read is a long process, but what parents do at home can really support emergent readers on the journey. | http://buddingreader.com/2015/09/ |
Hi everyone! I hope everyone had a wonderful long Memorial Day weekend! I did! I got to spend time with my son, daughter in law and grandson…I even got to watch Baby O while his mom and dad went out for the evening…that was fun. 😊
Well, I would like to share my latest project with you all. Since I am in a swap on instagram I thought I would combined the swap with my project I created for Eileen Hull this month.
I started by applying Eileen’s newest product, Easy Cut Adhesive (Eileen’s etsy shop is out at the moment but I went a head and tagged her shop which can be found HERE) to the Chipboard and then I picked cardstock for both the outside and inside covers from Crate Paper Chasing Dreams Paper Pad. I peeled the backing from the adhesive sheet and placed the cardstock on the cover. I did the same for the inside cover too. I then placed on top of the Notebook Diecut and ran through the Diecut Machine. I cut the notebook in half at the “center fold” and then folded both of the small “flaps” down. I cut a two inch wide by five inch long piece of chipboard and then adhered the small flaps to either side of the 2″x5″ of chipboard. I wanted to add more journals to the notebook and this was the perfect solution to extending the spine.
I cut and adhered a smaller piece of chipboard to fit between the flaps on the spine.
I measured in 1/2 inch increments and marked the holes at the top and bottom of the spine.
I punched the holes in the top and bottom of the spine with a Crop-A-Dile. I cut cardstock and Easy Cut Adhesive to fit the spine. I applied the adhesive sheet to the cardstock and then placed on the spine. I punched over the holes again.
I decorated the cover with some past Crate Paper chipboard/embellishemnt pieces and fussy cut butterflies and leaves from one of the paper pads (retired).
Before adding elastic cording I added eyelets to the holes of the spine. I also added a mini envelope to the back cover. I only glued two of the sides down so notes and tags could be added behind the envelope.
I created a folder from a file folder and used an assortment of cardstock and stickers (Crate Paper) to decorate the cover.
I adhered pom pom trim to the inside edge and added sticky notes just in case my swap partner needs to jot down something.
This little notebook has blank pages in it and I thought it would be perfect to either draw or even do memory keeping in it. I decorated the cover with more cardstock and stickers.
The last two notebooks I bought in a pack of three at the Dollar Tree (you can find these notebooks in packs of two at Walmart). The size of these notebooks are 3.25″x4.5″.
Like the other notebook and folder I decorated both sides with patterned cardstock, stickers, trim and chipboard/diecuts.
I slid all the notebooks and folder underneath the elastic cording.
I love how the pom pom trim looks from the side of the notebook planner! 😍 Also be sure to hop on over to Eileen Hull’s blog and see what the Inspiration Team has been up to! You can find the posts for this month HERE and HERE.
As always if you have any questions about this project please let me know.
I would also like to submit this project for the following challenge:
Word Art Wednesday: Anything Goes #388-#389
SUPPLIES:
Sizzix/Eileen Hull: Notebook Diecut
Sizzix: Diecut Machine
Hole Punch/Crop-A-Dile
Beacon Adhesive: Zip Dry Adhesive
Thread
Gems
Pom Pom Trim
Ric RacTrim
(Products listed above are Affiliated Links for your convenience)
Thanks for stopping by, hope your Tuesday is wonderful and enjoy creating! | https://creativepaperaddiction.com/2019/05/28/wings-and-things-notebook-planner/ |
Q:
Measuring relativity of simultaneity. Explanation of the bomb "paradox"?
Ok, so I know there have been some variations of this passed around and answered already, but I still can't quite understand how this works, so I want to clarify some particular points in this.
Lets start with Einstein's relativity of simultaneity thought experiment. Specifically the variant where a person stands in the middle of a train car with a lightbulb and another one is stationary on the ground next to the train as it passes by. When the moving observer lights the lightbulb he will wee the light reach both ends of the rail car at the same time. Meanwhile the stationary observer will see the light reach the rear end of the train car first.
I get the experiment itself. However I can't quite understand how it works if we start trying to measure the difference in relativity.
Say we put a light detector and an attosecond-accurate clock in each end of the train car and sync the clocks before the experiment (when both observers start in the same "stationary" frame of reference and can agree that the clocks a synced). When the light detector registers a photon, the clocks save the current timestamp and send it to a central computer. When the central computer receives both timestamps it compares them and outputs "EQUAL" or "DIFFERENT" on a screen.
In the original thought experiment we talk of the observers "observing" all the events take place, so, using the same terminology, the observer in the train car will, obviously see the light hit the detectors at the same time. The detectors will both be showing the same time (both will be showing t0 for example) when the light reaches both of them. He will, therefore, see the detectors record identical timestamps at that moment and send them to the central computer.
The stationary observer will observe the light hit the detectors at different moments in time. As both clocks are in the same frame of reference there is no time dilation between them, and the stationary observer will see them showing the same time at any specific point in time (even if they are slightly diverged from a similar clock he might have). Therefore he will see them both showing, say, t1 when light hits the first one, and both showing t2 when light hits the second one, which is fundamentally different from the first case. He will see them record different timestamps and send them to the central computer.
So now we have one person who saw the computer receive 2 identical timestamps, and another who saw it receive 2 different timestamps. So they will observe the computer perform the calculations and output different answers.
I understand there must be an error in such logic, but I can't understand where and why. Theoretically both observers could see the whole process happening up to and including the displaying of the result on the screen. At the same time, if the train stops and both observers walk up to the screen you would expect them to agree as to what is output on it.
Basically the specific questions I want to understand are:
If you are the stationary observer, then you should theoretically be able to observe the whole process and see "DIFFERENT" on the screen. Or if not, then how could you see anything else?
If you are the moving observer, then you should theoretically be able to observe the whole process and see "EQUAL" on the screen. Or if not, then how could you see anything else?
If anyone can explain this I would be extremely grateful, because this is frying my brain.
A:
The clocks are synched in the platform frame, so can't be synched in the train frame.
The story in the platform frame: The light beam takes longer to get to one of the synched clocks than to the other. Therefore it hits one clock when it says 1PM and the other when it says 2PM.
The story in the train frame: The light beam hits both clocks at the same time. However, one clock runs an hour behind the other, so the readings on the clocks when it hits are 1PM and 2PM.
(Or, if the clocks are synched in the train frame --- and therefore not in the platform frame --- you can tell essentially the same story in reverse.)
Edited to add in response to the OP's comment:
Why the clocks can't be synchronized in both frames: Because two points determine a line. If there is a frame in which both both clocks read 12:00 at the same time, then the line connecting those two events is a line of simultaneity for that frame. With just one spatial dimension, this is enough to uniquely determine the frame.
How the clocks get out of synch: You haven't given us enough information to answer this question. It depends on how the train starts moving.
Scenario A: In the platform frame, all parts of the train suddenly start moving rightward at the same time. Then in the (final) train frame, the train is initially moving left (so its synchronized clocks were both running slow) but then it stops. Moreover, the left side of the train stops before the right side does. Therefore there's a period when the right clock is running slow and the left clock is running normally. Therefore they get out of synch. (And incidentally, during the time when the right side is moving but the left side isn't, the train stretches.)
Scenario B: In the (final) train frame, all parts of the train suddenly start moving rightward at the same time. Then in the platform frame, the train is initially still (so its synchronized clocks run normally) but eventually starts moving. Moreover, the left side of the train starts moving before the right side does. Therefore there's a period in which the left clock runs slow while the right clock runs normally. Therefore they get out of synch. (And incidentally, during the period when the left side is moving but the right side isn't, the train shrinks.)
How I knew all this: I drew the spacetime diagram, which is the best way to solve almost any problem in relativity.
| |
By Taleh Mursagulov - Trend:
Azneft production union of Azerbaijan’s state oil company SOCAR has announced a tender to buy industrial furnace.
The participation fee is 354 manats.
Those who wish to participate in the tender should submit an application before 17:30 on December 19, 2018, and the tender offer before 11:30 on December 26, 2018. Opening of the tender packages will take place on December 26 at 11:30.
For more information please contact:
Address: AZ1000, Baku, Sabail District, Neftchilar Avenue, 73;
Tel.: (+994 12) 521-12-62 / (+994 12) 521-11-72;
In-company mobile: (+994 50) 841-12-62;
Email: [email protected]
Contact person: Mirismayil Orujov. | https://en.trend.az/business/tenders/2987927.html |
In this science experiment, milk is added to the water so that the light from the flashlight is more easily seen because the particles in the milk reflect the light. When the milk / water solution is viewed from the top, it appears to be a light blue colour because the particles separated out the blue waves of light.
Similarly, the earth’s atmosphere contains tiny drops of water and small particles of dust which causes the light from the sun to bend, resulting in the familiar blue skies! An effect called ‘Rayleigh scattering’ is when the molecules contained in the atmosphere scatter the different wavelengths of light coming from the sun. Abundant molecules in the atmosphere such as nitrogen and oxygen reflect light waves with shorter wavelengths, like violet and blue, much better than the longer wavelengths like red or orange. Wherever you look in the sky during the day, you can see scattered blue light overhead. | https://www.experiland.com/science-projects/find-out-why-the-sky-is-blue-in-this-science-experiment/ |
I knew it was BS. Barely drank any water specifically, and still physically healthy enough to do kickboxing.
If I only drank 30-50 ounces a day, I would be severely dehydrated. I guess most people don't move enough to sweat.
The eight-glasses-of-water misunderstanding comes from a WWII era army medical study of how much total moisture a cohort of healthy soldiers was consuming. This is pretty well known, so this article would appear to be researched poorly. The amount referenced was the volume of ALL MOISTURE consumed, including moisture from foodstuffs. A volume in liters was determined. Years later some idiot translated this total volume into "glasses", a measure which he thought would be more understandable to the general (American) public. And consequently, many people came to understand that one should actually drink eight glasses of water. But, as the article states, we get half that volume "automatically" from moisture in our food, and the remaining volume should be understood to include all liquids consumed, such as milk or soda. | https://www.newser.com/story/210219/drink-8-glasses-of-water-a-day-not-so-fast.html |
Share and Cite:
Received 18 February 2016; accepted 22 July 2016; published 25 July 2016
1. Introduction
For the online scheduling on a system of m uniform parallel machines, denoted by, each machine has a speed, i.e., the time used for finishing a job with size p on is. Without loss of generality, we assume. Cho and Sahni are the first to consider the on-line scheduling problem on m uniform machines. They showed that the LS algorithm for has competitive ratio not greater than. When and, they
proved that the algorithm LS has a competitive ratio and the bound is achieved when.
For, Epstein et al. showed that LS has the competitive ratio and is an optimal online algorithm, where the speed ratio.
Cai and Yang considered. Let and be two speed ratios. They showed that the algorithm LS is an optimal online algorithm when the speed ratios, where
and. Moreover, for the general speed ratios, they also presented an upper bound of the competitive ratio.
Aspnes et al. are the first to try to design algorithms better than LS for. They presented a new algorithm that achieves the competitive ratio of 8 for the deterministic version, and 5.436 for its randomized variant. Later the previous competitive ratios are improved to 5.828 and 4.311, respectively, by Berman et al. .
Li and Shi proved that for LS is optimal when and and presented an online algorithm with a better competitive ratio than LS for. Besides, they also showed that the
bound could be improved when and. For and, Cheng et al. proposed an algorithm with a competitive ratio not greater than 2.45.
A generalization of the Graham’s classical on-line scheduling problem on m identical machines was proposed by Li and Huang - . They describe the requests of all jobs in terms of order. For an order of the job, the scheduler is informed of a 2-tuple, where and represent the release time and the processing time of the job, respectively. The orders of request have no release time but appear on-line one by one at the very beginning time of the system. Whenever the request of an order is made, the scheduler has to assign a machine and a processing slot for it irrevocably without knowledge of any information of future job orders. In this on-line situation, the jobs’ release times are assumed to be arbitrary.
Our task is to allocate a sequence of jobs to the machines in an on-line fashion, while minimizing the maximum completion time of the machines. In the following of this paper, m parallel uniform machines which have speeds of respectively, are given. Let be any list of jobs, where job is given as order with the information of a release time and a processing size of.
The rest of the paper is organized as follows. In Section 2, some definitions are given. In Section 3, an algorithm U is addressed and its competitive ratio is analyzed.
2. Some Definitions
In this section we will give some definitions.
Definition 1. We have m parallel machines with speeds. Let be any list of jobs, where jobs arrives online one by one and each has a release time and a processing size of. Algorithm A is a heuristic algorithm. Let and be the makespan of algorithm A and the makespan of an optimal off-line algorithm respectively. We refer to
as the competitive ratio of algorithm A.
Definition 2. Suppose that is the current job to be scheduled with release time and size. We say that machine has an idle time interval for job, if there exists a time interval satisfying the following two conditions:
1) Machine is idle in interval and a job with release time is assigned on machine to start at time.
2).
It is obvious that if machine has an idle time interval for job, then we can assign to machine in the idle interval.
In the following we consider m parallel uniform machines with speeds and a job list with information for each job, where and represent its release time and size, respectively. For convenience, we assume that the sequence of machine speeds is non-decreasing, i.e.,
3. Algorithm U and Its Performance
Now we present the algorithm U by use of the notations given in the former section in the following:
Algorithm U:
Step 0. (*start the first phase*)
, ,.
Step 1. If there is a new job with release time and processing size given to the algorithm then go to Step 2. Otherwise stop.
Step 2. If there is a machine which has an idle time interval for job, then we assign to machine in the idle interval. Set and go to Step 1.
Step 3. Set. If then set, , Go to Step 1.
Step 4. (*start a new phase*)
Set, , and go to Step 3.
Now we begin to analyze the performance of algorithm U.
The following statement is obvious:
Lemma 1. Let be the stream of jobs scheduled in phase h and is the first job assigned in phase. Let be the largest load in an optimal schedule for job list. Then we have
Proof: If, let r be the fastest machine whose load does not exceed, i.e. . If there is no such machine, we set. If, then. It is
obvious that Hence we have
It means that can be assigned to the fastest machine in phase h. It is a contradiction to the fact that is the first job assigned in phase. Define, the set of machines with finishing time bigger than by the end of phase h. Since,. Denote by and the set of jobs assigned to machine by the on-line and the off-line algorithms, respectively. Since for any job the following inequalities hold
we get:
That means:
This implies that there exists a job () such that, i.e. there exists a job assigned by the on-line algorithm to a machine and assigned by the off-line algorithm to a slower machine.
By our assumptions, we have. Since, machine is at least as fast as machine, and thus. Since job was assigned before job and, we have
This implies
But this means that the on-line algorithm should have placed job on or a slower machine instead of, which is a contradiction. ¢
Theorem 2. Algorithm achieves a competitive ratio of 12.
Proof: Let denote the maximum load generated by jobs that were assigned during phase h; denote the last phase by. By the rules of our algorithm we have and
Hence the total height generated by the assignment is:
The claim of the theorem is trivially true if. For, phase h is started only if. In particular we have
Therefore we have
¢
4. Concluding Remarks
In this paper, we consider on-line scheduling for jobs with arbitrary release times on uniform machines. An algorithm with the competitive ratio of 12 is given. It should be pointed out that more detailed consideration should be taken in order to improve the competitive ratio.
Acknowledgements
The authors would like to express their thanks to the National Natural Science Foundation of China for financially supporting under Grant No. 11471110 and No. 61271264.
NOTES
*Corresponding author.
Conflicts of Interest
The authors declare no conflicts of interest.
References
●
Open Special Issues
●
Published Special Issues
●
Special Issues Guideline
Copyright © 2020 by authors and Scientific Research Publishing Inc.
This work and the related PDF file are licensed under a Creative Commons Attribution 4.0 International License. | https://www.scirp.org/journal/paperinformation.aspx?paperid=68954 |
You will spend time with a gentleman with a learning disability, visiting him 6 times a year.
He is in his late 50's and lives in a care home in Teddington.
He likes baking cakes, arts & crafts and going out for a little walk. He also has an interest in sports cars and trains.
You will spend time getting to know him and joining in with his favourite activities.
He will enjoy spending time with a friendly volunteer like you. By visiting every other month, you will make a huge difference.
You will be introduced and get to know him with the support of Mencap staff.
The times required can be flexible to your availability. Visits are every 2 months.
Visits will last around 1-1.5 hours.
You will make 6 visits per year, and we ask for a minimum commitment of 1 year.
This person lives in Teddington. Before applying please check your journey to make sure it's practical for you.
As a Mencap Visiting Volunteer you will be:
Mencap is the leading learning disability charity in England, Wales and Northern Ireland. We work with people with a learning disability and their families to challenge prejudice and change laws, and we directly support thousands of people to live their lives as they choose.
We have an ambitious vision for the UK to be the best place in the world for people with a learning disability to live happy and healthy lives.
Volunteering with us is YOUR opportunity to help us achieve this, whilst having the chance to develop your skills, meet new people and join a passionate and dedicated team. | https://jobs.mencap.org.uk/vacancies/25152/mencap-volunteer-visitor--teddington.html |
80205 Market Overview
There are 7 homes currently listed for 80205. 80205 real estate market trends show that home prices range from $614,900 to $975,000, and the median sales price in the 80205 zip code is $682,550. The total number of properties sold within the past twelve months is 868. In 80205, 2 properties are in foreclosure, 0 are bank owned properties, and 2 are headed for auction.
Active Residential Inventory
0
Total Inventory
Recently Sold Properties
80205 Home Values
$500,000
The median list price of a home in Denver, is $500,000. A total of 7 properties are for sale, and the percentage of properties for sale in the zip code is 77.8%. The total count of 80205 properties is 9,833.
80205 Real Estate Market Overview
80205 Foreclosures Market Overview
Last 24 Hours
There is 1 city within 80205; within this city, the median estimated home value for homes in foreclosure ranges from $1,606,484 in Denver to the lowest value of $198,900 in Denver. Foreclosure homes account for 0.02% of properties in 80205 with Denver containing 126 foreclosures, the highest number of foreclosure properties in a single city. | https://www.realtytrac.com/market-trends/80205/ |
Round-the-world solar plane takes off from Nanjing for Hawaii
NANJING, China – The first round-the-world solar-powered flight continues its journey, taking off from Nanjing, China, to cross the Pacific Ocean towards Hawaii.
The flight took off from China on Sunday (May 31), as it continues its 35,000 km (22,000 mile) journey seeking to demonstrate that flying long distances fueled by renewable energy is possible.
The Solar Impulse 2 took off from Nanjing in the early hours of Sunday to cross the Pacific Ocean for Hawaii, before flying across the United States and southern Europe to arrive back in Abu Dhabi by late July.
Pilots Andre Borschberg and Bertrand Piccard will take turns at the controls of Solar Impulse 2, which began its journey in Abu Dhabi in the United Arab Emirates in March, as it makes its way around the globe in about 25 flight days at speeds of between 50 kph and 100 kph (30 mph to 60 mph).
The aircraft is as heavy as a family car at 2,300 kg (5,100 lb) but has a wingspan as wide as the largest airliner.
The design and construction of the Solar Impulse took 12 years. | https://92newshd.tv/round-the-world-solar-plane-takes-off-from-nanjing-for-hawaii/ |
Sometimes you have many records or instances in a collection and you want to limit the amount of data returned by a query.
Strategy
Limiting Records in a Document
If you are limiting records in a large XML document you can do this by adding a predicate to the end of your for loop:
for $person in doc($file)/Person[position() lt 10]
Using Subsequence to Limit Results
The following query retrieves only the first 10 documents, in sequential document order, contained within a collection.
for $person in subsequence(collection($my-collection)/person, 1, 10)
Where:
subsequence($sequence, $starting-item, $number-of-items)
Note that the second argument is the item to start at and the third is the total number of items to return. It is NOT the last item to return.
Sorting Before Limiting Results
Note that usually you will want to get the first N documents based on some sorting criteria, for example people that have a last name that start with the letter "A". So the first thing you must do is create a list of nodes that have the correct order and then get the first N records from that list. This can usually be done by creating a temporary sequence of sorted items as a separate FLWOR expression.
let $sorted-people := for $person in collection($collection)/person order by $person/last-name/text() return $person for $person at $count in subsequence($sorted-people, $start, $number-of-records) return <li>$person/last-name/text()</li>
Adding Buttons to Get Next N Records
Getting the Next N Rows
After you fetch the first N items from your sequence you frequently want to get the next N rows. To do this you will need to add buttons to your report for "Previous N Records" and "Next N Records". Those buttons will pass parameters to your XQuery telling where to start and how many records to fetch.
To do this you will need to call your own script with different parameters on the URL. To keep track of the URL you can use the get-url() function that comes with eXist. For example:
let $query-base := request:get-url()
.
If your query was run from http:/localhost:8080/exist/rest/db/apps/user-manager/views/list-people.xq this is the string that would be returned.
We will also get two parameters from the URL for the record to start at and the number of records to fetch:
let $start := xs:integer(request:get-parameter("start", "1")) let $num := xs:integer(request:get-parameter("num", "20"))
Now we will create two HTML buttons that allow us to get next N records or the previous N records. | https://en.wikibooks.org/wiki/XQuery/Limiting_Result_Sets |
proof of PTAH inequality
In order to prove the PTAH inequality two lemmas are needed. The first lemma is quite general and does not depend on the specific and that are defined for the PTAH inequality.
The setup for the first lemma is as follows:
We still have a measure space with measure . We have a subset . And we have a function which is positive and is integrable in for all . Also, is integrable in for each pair .
Define by
and
by
Lemma 1 (1)
(2) if then . If equality holds then a.e [m].
Proof It is clear that (2) follows from (1), so we only need to prove (1). Define a measure . Then
so we can use Jensen’s inequality for the logarithm.
The next lemma uses the notation of the parent entry.
Lemma 2 Suppose for and . If then
Proof. Let . By the concavity of the function we have
where for .
so that
|(1)|
It is enough to prove the lemma for the case where for all . We can also assume for all , otherwise the result is trivial.
Let and so that .
Raise each side of (1) to the power:
|(2)|
so that
|(3)|
Multiply (3) by to get:
|(4)|
Claim: There exist , such that
|(5)|
If so, then substituting into (4)
So it remains to prove the claim. We have to solve the system of equations , for . Rewriting this in matrix form, let , , and , where and if , . The columns sums of are , since . Hence is singular and the homogenous system has a nonzero solution, say . Since is nonsingular, it follows that . It follows that for some and therefore . If necessary, we can replace by so that . From (5) it follows that for all .
Now we can prove the PTAH inequality. Let .
We calculate by differentiating under the integral sign. If then
Thus
|(6)|
If then by writing
where it is clear that each integral is 0, so that . So again, (6) holds. Therefore,
Then
Now by Lemma 1, with we get . | https://planetmath.org/ProofOfPTAHInequality |
ROME: Spanish climber Mikel Landa crowned a marathon six-hour slog with a majestic late charge to clinch victory at the summit finish on stage four of the Tirreno-Adriatico on Saturday.
Recently signed with Movistar, the former Sky man Landa timed his run late and finished the day in a strong position to challenge for the overall victory, trailing new blue jersey holder Damiano Caruso by 20 seconds.
The stage was marked by overnight leader Geraint Thomas of Sky suffering an untimely mechanical problem and having to change bikes on the final climb, he lost 40 seconds on the day while his muted team leader Chris Froome was 1min 10sec off the pace on Saturday too.
The race also suffered the loss of Giro d’Italia winner and world champion Tom Dumoulin, who fell and pulled out early and was later reported to have suffered only cuts and bruises.
For the pure climbers, Dumoulin’s demise takes much of the fear out of Tuesday’s final-day time-trial and the effects that would have on the overall standings.
Landa took 6hrs, 22mins and 13secs to complete the challenging stage by outsprinting Rafal Majka and George Bennett to the line. | https://www.pressreader.com/malaysia/the-borneo-post-sabah/20180312/281891593791841 |
BACKGROUND OF THE INVENTION
SUMMARY OF THE INVENTION
DETAILED DESCRIPTION OF THE INVENTION
EXAMPLES 1 TO 11 AND COMPARATIVE EXAMPLES 1 AND 2
(1) Moment of Inertia
(2) Flight Distance and Flight Directionality
The present invention relates to a wood-type golf club head having stable performances of flight direction and flight distance of hit ball.
FIG. 8A
FIG. 8A
As shown in as to a wood-type golf club head “a” for right-handed golfers (all explanations made herein being for right-handed golfers), if a golf ball “b” is struck by the wood-type golf club head “a” at a position on a toe side of a sweet spot SS that is a point at which a normal line drawn from the center of gravity G of the club head with respect to a face “f” intersects the face “f” (such a hitting may be hereinafter referred to as “toe hit”), the club head rotates clockwise about the center of gravity G. In the ball which is in contact with the face, a side spin which causes the ball to rotate in the counterclockwise direction opposite to the rotation of the club head (i.e., so-called hook spin for right-handed golfers) generates. On the other hand, if the club head hits the ball on a heel side of the sweet spot SS (such a hitting may be hereinafter referred to as “heel hit”), the club head rotates counterclockwise about the center of gravity G, and the ball which is in contact with the face causes a side spin of the clockwise rotation (so-called slice spin). Such a phenomenon is well known as a gear effect. If the face “f” is flat, as shown in , the ball “b” is struck out approximately parallel to the target flight direction “j” and thereafter the ball is driven to the left or right direction to cause a hook or slice.
FIG. 8B
In order to solve such a problem of poor flight direction performance, a face bulge has been conventionally provided to the face “f” of the wood-type golf club head “a”, as shown in . The face bulge is a rounded or curved surface having a radius of curvature Rx which is smoothly and slightly convex toward the front side. In case of a toe hit by a club head provided with a face bulge, the ball is struck out at a deflection angle θ with respect to the target flight direction “j” by the face bulge, and then would hook by a side spin to curve back toward the target flight direction. In case of a heel hit, the ball is struck out at a deflection angle to the left of the target flight direction “j”, and then would slice by a side spin to curve back toward the target flight direction. Therefore, the face bulge serves to improve the flight direction performance by a side spin caused by a horizontal gear effect.
FIG. 9A
1
2
A similar phenomenon to the horizontal gear effect also takes place about a horizontal axis passing through the head's center of gravity G in the toe-heel direction. This may also be called “vertical gear effect”. For example, as shown in , if the club head “a” hits the golf ball “b” on a crown “c” side above the sweet spot SS of the face “f” (such a hitting may be hereinafter referred to as “high hit”) or on a sole “s” side below the sweet spot SS (such a hitting may be hereinafter referred to as “low hit”), the club head rotates about the horizontal toe-heel axis by a moment which is the product of a force F received from the golf ball “b” and a vertical distance L or L between a hitting position and the center of gravity G. At that time, the golf ball contacting the face “f” receives a force acting in the direction opposite to the rotation of the club head “a” by a frictional force. Thus, off-center hits above the sweet spot reduce the amount of backspin of the golf ball “b”, and off-center hits below the sweet spot increase the amount of the backspin.
FIG. 9A
In order to prevent such an uneven backspin amount from occurring, a face roll has been conventionally provided to the face “f” of the wood-type golf club head “a”, as shown in . The face roll is a rounded or curved surface having a radius of curvature Ry which is smoothly and slightly convex toward the front side when the face is viewed from the side.
FIGS. 9A and 9B
FIG. 9B
FIG. 9A
1
2
In case of hitting a golf ball “b” by club heads “a” shown in at the same position below or above the sweet spot SS, the vertical distance L or L between the force F and the center of gravity G is larger when hitting by a face “fn” having no face roll of the club head shown in as compared with the hitting by a face “f” having a face roll of the club head shown in . In other words, the rotational moment is reduced as the radius of curvature Ry of the face roll is reduced and, therefore, the vertical gear effect can be reduced. Further, in case of off-center hits above the sweet spot, the face roll increases the launch angle δ. This is useful for preventing the flight height from lacking with decrease in the amount of backspin on off-center hits above the sweet spot. The face roll reduces the launch angle δ for off-center hits below the sweet spot and, therefore, it is also effective in preventing a hit ball from flying too high owing to increased amount of backspin.
2
FIG. 8C
Size increase of golf club heads has progressed rapidly with recent development of thin wall molding technology for metal materials. Large-sized golf club heads enable to have a large moment of inertia about a vertical axis passing through the center of gravity. For example, club heads having a moment of inertia of 4,000 g·cmor more are known. However, in case of club heads having a large moment of inertia about the vertical axis, the amount of rotation or twisting of the head “a” about the vertical axis is small for both the toe hit and the heel hit. Thus, the amount of sidespin imparted to the ball by the horizontal gear effect is also small. Therefore, as shown in , club heads having, for example, a small face bulge radius Rx (i.e., a large curvature) and a large moment of inertia about the vertical axis cannot impart a sidespin in an amount commensurate with the deflection angle θ of the hit ball, in the toe hits or the heel hits, despite that the deflection angle θ becomes large. Since the hit ball does not curve back to the target flight direction “j” by such a reason, large-sized club heads have a problem of poor flight direction performance.
1
2
1
1
2
JP-A-2001-161866 discloses a golf club head having a horizontal bulge radius R of 480 to 765 mm and a vertical roll radius R larger than the radius R. JP-A-8-089603 disclose a golf club head having a horizontal bulge radius R of at most 9 inches and a vertical roll radius R of at most 9 inches. In these prior art, the horizontal bulge radius and the vertical roll radius are not specified in association with the moment of inertia about the vertical axis of the head. Thus, these proposed golf club heads still have room for improvement.
It is an object of the present invention to provide a wood-type golf club head having an improved directional stability of hit ball and an increased flight distance.
This and other objects of the present invention will become apparent from the description hereinafter.
2
It has been found that the flight directionality and the flight distance of wood-type golf club heads can be improved when, with respect to club heads having a moment of inertia about the vertical axis passing through the club head's center of gravity as large as 4,000 to 5,900 g·cm, the radius of curvature Rx of the face bulge is restricted to a specific range and the radius of curvature Ry of the face roll is determined in association with the bulge radius Rx.
In accordance with the present invention, there is provided a hollow wood-type golf club head having a face for hitting a golf ball, wherein:
2
the moment of inertia Ix about a vertical axis passing through the club head's center of gravity is 4,000 to 5,900 g·cmin the standard state that the club head is placed on a horizontal plane in the state that an axial center line of a shaft is disposed in an optional vertical plane and is inclined at a prescribed lie angle with respect to the horizontal plane and said face is inclined at a prescribed loft angle,
said face has a face bulge and a face roll,
the radius of curvature Rx of the face bulge is from 12 to 25 inches (about 30.48 to about 63.50 cm), and
the Ry/Rx ratio of the radius of curvature Ry of the face roll to said radius of curvature Rx of the face bulge is 0.50 to 0.90.
2
It is preferable that the Rx/Ix ratio of the face bulge radius Rx to the moment of inertia Ix is from 0.0030 to 0.0045 inch/g·cm.
2
It is also preferable that the Ry/Iy ratio of the face roll radius Ry to a moment of inertia Iy about a horizontal axis extending through the club head's center of gravity in the toe-heel direction is from 0.0030 to 0.0080 inch/g·cm.
Further, it is preferable that the golf club head satisfies the following relationship:
Ry/Rx
Iy/Ix
1.0≦()/()≦2.0
2
2
wherein Rx is the radius of curvature of the face bulge (inch), Ry is the radius of curvature of the face roll (inch), Ix is the moment of inertia about the vertical axis (g·cm), and Iy is the moment of inertia about the horizontal axis (g·cm).
2
The wood-type golf club heads of the present invention have a large moment of inertia Ix about a vertical axis extending through the club head's center of gravity, i.e., 4,000 to 5,900 g·cm, and a relatively large radius of curvature Rx of the face bulge, i.e., 12 to 25 inches. In case of golf club heads having such large moment of inertia and large face bulge radius, the twisting of the club heads is small for off-center hits on the toe or heel side, so occurrence of the horizontal gear effect is suppressed. However, since the golf club heads of the present invention have a face bulge having a large radius of curvature, the deflection angle of hit ball is small for the toe or heel side shots. Therefore, golf balls are struck out with a reduced amount of sidespin at a small deflection angle with respect to the target flight direction, and mildly curve back toward the target. Therefore, the golf club heads of the present invention have a stabilized flight directionality and an increased flight distance performance.
On the other hand, if the radius of curvature of the face bulge is large, the hitting face becomes easy to bend by impact of the ball since the face is flattened. Therefore, such a face has a possibility of imparting an excessively high spring effect (rebound property) to golf club heads, thus resulting in violation of golf rules. In contrast, in the club heads of the present invention, the radius of curvature Ry of the face roll is made smaller than the radius of curvature Rx of the face bulge so that the Ry/Rx ratio falls within the range of 0.50 to 0.90. The flexural rigidity of a face portion of the heads can be increased by a large curvature of the face roll. Therefore, the spring effect of the club heads can be prevented from excessively increasing without adopting any other means, for example, without increasing the wall thickness of the face portion.
An embodiment of the present invention will be explained below with reference to the accompanying drawings.
FIGS. 1 to 4
FIG. 3
1
are perspective, front and plan views of a wood-type golf club head according to an embodiment of the present invention, and an enlarged cross sectional view along the line A-A of , respectively.
1
1
1
1
a
b
a
The wood-type golf club head in this embodiment comprises a head body having a hollow structure, and a hosel portion which is disposed on a heel side of the body for inserting a shaft.
1
3
2
4
2
2
1
5
2
2
1
6
4
5
2
2
2
2
1
1
a
a
b
c
d
a
The head body includes a face portion having a face for hitting a golf ball on its front side, a crown portion which extends from an upper edge of the face and forming the upper surface of the head , a sole portion which extends from an lower edge of the face and forming the bottom surface of the head , and a side portion which extends between the crown portion and the sole portion to connect them from a toe side edge of the face to a heel side edge of the face through a back face BF of the head . The head body has a hollow portion “i”.
1
4
1
7
7
7
1
b
a
The hosel portion is disposed on a heel side of the crown portion of the head body and has a cylindrical shaft inserting hole to attach a shaft (not shown). Since the axial center line of the shaft inserting hole substantially agrees with the axial center line CL of the shaft when the shaft is inserted into the hole , it is used as the axial center line CL of the shaft when no shaft is attached to the club head .
FIGS. 1 to 4
FIGS. 2 to 4
1
1
2
1
1
1
In , the club head is kept in the standard state. The term “standard state” of a golf club head as used herein denotes the state that, as shown in , golf club head is placed on a horizontal plane HP in the state that the axial center line CL of a shaft is disposed in an optional vertical plane VP and is inclined at a prescribed lie angle β with respect to the horizontal plane HP, and the hitting face is inclined at a prescribed loft angle α (real loft angle, hereinafter the same) given to the head . The head referred to herein is in the standard state unless otherwise noted. The terms “prescribed lie angle β” and “prescribed loft angle α” as used herein denote those previously given to the head .
1
1
1
1
2
2
1
1
2
1
2
FIG. 3
FIG. 3
Further, with respect to the club head , the up-down direction and the height direction denote those of the club head in the standard state. The front-rear direction denotes, when the head in the standard state is viewed from above, namely in a plane view of the head (), a direction which is parallel to a perpendicular line N drawn from the club head's center of gravity G to the face , in other words, a direction parallel to a line N connecting the center of gravity G and a sweet spot SS. A face side is the front and a back face BF side is the rear or back. The toe-heel direction of the club head denotes a direction parallel to the vertical plane VP defined above, in other words, a direction perpendicular to the front-rear direction, in the plan view of the head in the standard state (). The sweet spot SS is a point where a normal line N drawn to the face from the center of gravity G of the head intersects the face .
The term “wood-type golf club head” does not mean that the head is made of a woody material, but means golf club heads having a so-called wood-type head shape, e.g., driver (#1 wood), brassy (#2 wood), spoon (#3 wood), baffy (#4 wood) and cleek (#5 wood), and comprehends heads which are different from these heads in number or name, but have a shape approximately similar to these heads.
1
1
1
The club head in this embodiment is produced from a metallic material. Preferable examples of the metallic material are, for instance, a stainless steel, a marageing steel, a pure titanium, a titanium alloy, an aluminum alloy, and combinations of these metals. For the purpose of weight reduction and so on, a nonmetallic material, e.g., fiber-reinforced resins and ionomers, may be used in a part of the club head . The club head can be produced by joining a plurality of members or pieces (e.g., two to five pieces). The number of pieces is not particularly limited. Each member or piece is formed by various molding methods, e.g., casting, forging and pressing.
1
1
7
1
1
1
3
3
3
3
3
Preferably, the club head of the present invention has a head volume of at least 400 cm, especially at least 420 cm, more especially at least 430 cm. The “head volume” denotes the volume of a portion surrounded by the outer surface of head whose shaft inserting hole in the hosel portion is filled up. Such a large head volume would provide a sense of ease to a player at the time of address, and it is also useful in improving the flight directionality since the moment of inertia or the depth of the center of gravity of the club head can be increased. On the other hand, if the volume of the club head is too large, problems may arise, e.g., increase of head weight, deterioration of swing balance, deterioration of durability and violation of golf rules. From such points of view, the volume of head is preferably at most 470 cm, more preferably at most 460 cm.
1
In consideration of swing balance, rebound property, swing easiness or the like, it is preferable that the weight of club head is at least 180 g, especially at least 183 g, more especially at least 185 g, and it is at most 220 g, especially at most 215 g, more especially at most 213 g.
1
1
1
2
2
2
2
The club head of the present invention has a moment of inertia Ix of 4,000 to 5,900 g·cmabout the vertical axis passing through the center of gravity G. The club head having such a large moment of inertia Ix about the vertical axis can diminish the amount of rotation (twisting) of the head about the vertical axis for toe or heel side off-center hits, whereby the amount of sidespin imparted to the ball by the horizontal gear effect is decreased to improve the straightness of flight. In order to more surely exhibit such an effect, the moment of inertia Ix about the vertical axis is preferably at least 4,100 g·cm, more preferably at least 4,200 g·cm. However, if the moment of inertia Ix is too large, it may violate golf rules which provide the upper limit of the moment of inertia Ix. Therefore, it is preferable that the moment of inertia Ix is at most 5,800 g·cm. Such a large moment of inertia can be easily realized by increasing the head volume to fall within the above-mentioned range, by adjusting the thickness of respective portions of the head, and/or by additionally using a material having a high specific gravity, so as to distribute the weight to the perimeter of the club head.
FIG. 5
FIG. 5
1
2
1
1
2
shows a horizontal cross section view passing through sweet spot SS of club head in the standard state. As apparent from the drawing, the face of the club head in this embodiment is provided with a face bulge FB which is smoothly convex toward the front of the head when the face is viewed from above. In this embodiment, the face bulge FB is provided to substantially the entire region of the face in the toe-heel direction. The convex curvature of this face bulge FB is formed not only at the section position shown in , but also smoothly extends upward and downward.
1
2
FIG. 8C
Since the club head in this embodiment has a large moment of inertia Ix about the vertical axis of 4,000 g·cmor more as stated above, the horizontal gear effect is reduced and, therefore, the amount of sidespin is also reduced. If a face bulge BF having a small radius of curvature Rx is provided to such a head, a golf ball is struck out at an excessive deflection angle θ on toe or heel hits (cf. ), and does not curve back toward the intended line of flight due to a reduced sidespin. Thus, in the present invention, it is required that the radius of curvature Rx of the face bulge FB is 12 inches (about 30.48 cm) or more. Like this, with respect to golf club heads having a large moment of inertia Ix about the vertical axis, by selecting a large value as a radius of curvature Rx of the face bulge FB in association with the moment of inertia Ix, the horizontal deflection angle θ of the hit ball on toe or heel hits is reduced to improve the flight direction performance. The radius of curvature Rx of the face bulge FB is preferably at least 13 inches, more preferably at least 14 inches.
1
On the other hand, though the club heads have a large moment of inertia Ix about the vertical axis, a slight horizontal gear effect still generates and it imparts a sidespin to the golf ball. Therefore, if the radius of curvature Rx of the face bulge FB is too large, the deflection angle θ on a mishit may become excessively small, so it cannot compensate for the gear effect spin and the flight directionality may be deteriorated. Therefore, from such a point of view, it is required that the radius of curvature Rx of the face bulge FB is at most 25 inches (about 63.50 cm). The radius of curvature Rx is preferably at most 24 inches, more preferably at most 22 inches.
FIG. 5
2
2
2
2
2
2
2
2
2
2
2
2
2
2
2
c
d
c
d
c
d.
c
d
Herein, as shown in , the “radius of curvature Rx of the face bulge FB” is defined, for convenience's sake, as a radius of a single arc which passes through the following three points; a face toe side point Pt on the face which is apart from the toe side edge of the face toward the sweet spot SS by a distance of 20 mm in the toe-heel direction, a face heel side point Ph on the face which is apart from the heel side edge of the face toward the sweet spot SS by a distance of 20 mm in the toe-heel direction, and the sweet spot SS. In the case that the face has score lines or the like, the radius of curvature Rx is defined for the face in the state that the score lines or the like are filled. Further, in the case that the toe side edge and the heel side edge of the face can be clearly identified by edge lines or the like, these edge positions are the edges and However, in the case that the toe side edge and the heel side edge are not clearly determined, they are defined as positions at which the actual radius of curvature “ra” of the face reaches 20 mm for the first time when the actual radius “ra” is measured from the sweet spot SS toward the toe and heel sides.
1
1
2
2
The larger the moment of inertia Ix about the vertical axis, the smaller the amount of rotation (twisting) of the club head about a vertical axis passing through the club head's center of gravity G on mishit, so that the amount of sidespin imparted to the ball by the horizontal gear effect is reduced. Therefore, it is necessary to diminish the deflection angle θ of a ball struck out by toe or heel hits in accordance with the reduction of sidespin. From such a point of view, it is effective to determine the radius of curvature Rx of the face bulge FB in association with the moment of inertia Ix of the club head about the vertical axis. Specifically, it is preferable that the ratio Rx/Ix of the radius of curvature Rx (inch) of the face bulge FB to the moment of inertia Ix (g·cm) about the vertical axis is at least 0.0028, especially at least 0.0030, more especially at least 0.0031, further especially at least 0.0033, still further especially at least 0.0037, and is at most 0.0050, especially at most 0.0045, more especially at most 0.0040. When the radius of curvature Rx of the face bulge is determined in association with the moment of inertia Ix so that the Rx/Ix ratio (inch/g·cm) falls within the above range, optimum combination of the amount of sidespin and the deflection angle θ of flight is realized, so the flight directionality and the flight distance can be further improved.
FIG. 4
FIG. 6
FIG. 4
FIG. 6
1
3
2
1
1
2
shows a vertical cross section view passing through the center of gravity G and the sweet spot SS of club head in the standard state. shows an enlarged partial view of the face portion in . As apparent from these drawings, the face of the club head in this embodiment is provided with a face roll FR having a radius of curvature Ry which is smoothly convex toward the front of the head when the face is viewed from the side. In this embodiment, the face roll FR is provided to substantially the entire region of the face from top to bottom. The convex curvature of this face roll FR is formed not only at the section position shown in , but also smoothly extends toward both the toe and heel sides.
In the present invention, the radius of curvature Ry of the face roll FR is set to 0.50 to 0.90 times the radius of curvature Rx of the face bulge FB. In other words, the radius of curvature Ry of the face roll is selected so that the Ry/Rx ratio falls within the range of 0.50 to 0.90.
3
3
3
3
3
If the radius of curvature Rx of the face bulge FB is increased, the face portion is more flattened, so the face becomes easy to bend by impact of a ball. Thus, a flattened face is likely to excessively enhance a spring effect (rebound property) of the club head, resulting in violation of golf rules. In the present invention, the flexural rigidity of the face portion is enhanced without increasing the thickness or the like of the face portion by setting the radius of curvature Ry of the face roll FR to a value smaller than the radius of curvature Rx of the face bulge FB, thereby preventing the spring effect from excessively appearing. The flexural rigidity of the face portion can be enhanced by increasing the thickness of the face portion , but the increase of the thickness is undesirable since the degree of freedom in weight distribution design is remarkably lowered.
FIG. 9A
If the Ry/Rx ratio is less than 0.50, the radius of curvature Ry of the face roll FR is excessively small as compared with that of the face bulge FB and, therefore, the face looks like it is protruding forwardly at the time of address and the sense of use is remarkably deteriorated. Further, as shown in , if the face is provided with a face roll, the launch angel δ on high or low hits is larger or smaller than the launch angle of a ball struck at sweet spot. When the radius of curvature Ry of the face roll FR is small, the difference in launch angle between the off-center hit and the center hit is noticeable and the flight distance tends to be unstable. From such points of view, the Ry/Rx ratio is preferably at least 0.55, more preferably at least 0.60.
3
1
On the other hand, if the Ry/Rx ratio is more than 0.90, the face roll FR is flattened, so an effect of enhancing the rigidity of the face portion is not sufficiently obtained. As a result, the rebound property of the club head tends to exceed an upper limit provided in golf rules, unless any measure such as thickening the face portion is taken. From such a point of view, the Ry/Rx ratio is preferably at most 0.85, more preferably at most 0.80.
FIG. 6
FIG. 6
2
2
2
2
2
2
2
9
9
2
2
2
2
2
2
2
2
a
b
a
b
a
b.
a
b
Herein, as shown in , the “radius of curvature Ry of the face roll FR” is defined, for convenience's sake, as a radius of a single arc which passes through the following three points; a face upper side point Pu on the face which is apart from the upper edge of the face toward the downside by a distance of 10 mm in the vertical direction, a face downside point Pd on the face which is apart from the lower edge of the face toward the upper side by a distance of 10 mm in the vertical direction, and the sweet spot SS. In the case that the face has score lines or the like, the radius of curvature Ry is defined for the face in the state that the score lines or the like are filled. Further, in the case that the upper edge and the lower edge of the face can be clearly identified by edge lines or the like, these edge positions are the edges and However, in the case that the upper edge and the lower edge are not clearly determined, they are defined as positions at which the actual radius of curvature “ra” of the face reaches 20 mm for the first time, in the vertical cross section as shown in , when the actual radius “ra” is measured from the sweet spot SS toward the upper and down sides.
The radius of curvature Ry of the face roll FR is not particularly limited so long as it satisfies the above-mentioned Ry/Rx ratio condition. The change in launch angle δ on high and low hits tends to become large as the radius of curvature Ry decreases. Therefore, it is preferable that the radius of curvature Ry is at least 8 inches, especially at least 9 inches, more especially at least 10 inches. On the other hand, the vertical gear effect on high and low hits appears more strongly as the radius of curvature Ry increases, so the amount of backspin tends to be not stabilized to result in unstable flight distance. From such points of view, it is preferable that the radius of curvature Ry is at most 20 inches, especially at most 19 inches, more especially at most 17 inches.
1
3
FIG. 9A
In the club head of the present invention, an excessive bending of the face portion by impact of a ball is suppressed to control appearance of the spring effect within provisions of golf rules by adopting a small radius of curvature for the face roll FR. However, when the radius of curvature of the face roll is small, there is a tendency that the launch angle δ for high hits becomes large and the launch angle δ for low hits becomes noticeably small (cf. ). In particular, loss of flight distance is easy to occur at the time of low hits. In order to prevent such a lowering of the flight distance, it is preferable that the moment of inertia Iy about a horizontal axis extending through the center of gravity G in the toe-heel direction is set to a small value to deliberately increase the amount of rotation of the head about the horizontal axis for mishits. Thus, a relatively strong vertical gear effect appears for low hits to impart an increased amount of backspin to a ball, whereby even if the launch angle δ is low, a large lift force is imparted to the ball to minimally suppress the lowering of the flight distance.
2
2
2
2
2
2
2
In order to obtain such an action, it is preferable that the moment of inertia Iy about the horizontal axis is at most 5,900 g·cm, especially at most 4,000 g·cm, more especially at most 3,600 g·cm. On the other hand, if the moment of inertia Iy about the horizontal axis is too small, the amount of rotation of the head about the horizontal axis becomes excessively large for mishits, so the flight distance tends to be not stabilized. Therefore, it is preferable that the moment of inertia Iy is at least 1,500 g·cm, especially at least 1,800 g·cm, more especially at least 2,000 g·cm, further especially at least 3,000 g·cm.
2
2
In particular, it is preferable that the Ry/Iy ratio (inch/g·cm) of the radius of curvature Ry (inch) of the face roll to the moment of inertia Iy (g·cm) about the horizontal axis is at least 0.0030, especially at least 0.0035, more especially at least 0.0040, and is at most 0.0080, especially at most 0.0070, more especially at most 0.0060.
Further, in order to stabilize the flight distance and the flight directionality, it is preferable that the Iy/Ix ratio of the moment of inertia Iy about the horizontal axis to the moment of inertia Ix of the vertical axis is at least 0.30, especially at least 0.35, more especially at least 0.40, and is at most 0.80, especially at most 0.75, more especially at most 0.70. If the Iy/Ix ratio is less than 0.30, the flight directionality is improved, but there is a tendency that the vertical gear effect appears strongly in excess for high and low hits, so the launch angle δ is not stabilized, thus resulting in unstable flight distance. On the other hand, if the Iy/Ix ratio is more than 0.80, the vertical gear effect on high and low hits is suppressed to stabilize the flight distance, but the flight directionality on toe and heel hits tends to deteriorate.
Further, as a result of inventor's investigation, it has been found that the flight distance and the flight directionality are further improved when the moment of inertia ratio Iy/Ix and the radius of curvature ratio Ry/Rx satisfy the following relationship:
Ry/Rx
Iy/Ix
1.0≦()/()≦2.0
If the ratio (Ry/Rx)/(Iy/Ix) is less than 1.0, the flight distance is easy to become unstable, and if it is more than 2.0, the flight directionality is easy to deteriorate.
While a preferable embodiment of the present invention has been described with reference to the drawings, it goes without saying that the present invention is not limited to only such an embodiment and various changes and modifications may be made.
The present invention is more specifically described and explained by means of the following examples. It is to be understood that the present invention is not limited to these examples.
FIGS. 1 to 4
Wood-type golf club heads having a base structure shown in were prepared according to the specifications shown in Table 1 and tested with respect to the flight distance and flight direction performances.
FIG. 1
Head weight: 200 g
3
Head volume: 460 cm
Real loft angle: 11°
Thickness of crown portion: 0.6 mm uniform thickness
FIG. 4
Height “h” of face: suitably varied within the range of 40 to 65 mm (cf. )
FIG. 5
Width FW of face: suitably varied within the range of 90 to 105 mm (cf. )
The respective club heads were prepared by welding a plate-like face member “k” and a hollow head body member “m”, the boundary of which is shown in by a chain line. Specifically, the head body member “m” was prepared by precision casting of a Ti-6Al-4V alloy. The face member “k” was prepared by subjecting an α-β titanium alloy (Titanium Alloy “SP700HM” made by JFE Steel Corporation having a composition of Al: 4.0 to 5.0% by weight, V: 2.5 to 3.5% by weight, Mo: 1.8 to 2.2% by weight, Fe: 1.7 to 2.3% by weight, and Ti and unavoidable impurities: the rest) to machine work and press work to have a thick center portion and a thin peripheral portion. The head body member “m” and the face member “k” were joined by laser welding to give club heads having the following common specifications.
FIG. 4
3
2
2
2
3
3
a
a
d,
b
c
The wall thickness of the head body member was partially changed so that the thickness of the side portion falls within the range of 0.5 to 1.5 mm and the thickness of the sole portion falls within the range of 0.7 to 2.0 mm. Further, the face member was formed into a periphery-thin wall structure (cf. ) having a center thick wall region with a similar shape to the shape of the face formed by the peripheral edges to a peripheral thin wall region and a transition portion between them, in which the difference in thickness between the thick wall region and the thin wall region was 1.0 mm, and the entire thickness was adjusted so that the CT value according to the Pendulum Test Protocol (test rules of the R & A) falls within the range of 250±20. Thus, the center of gravity and the moment of inertia were adjusted.
The testing methods are as follows:
The moment of inertia was measured using Moment of Inertia Measuring Instrument Model No. 005-004 made by INERTIA DYNAMICS INC.
Same FRP shafts were attached to all club heads to be tested to give wood-type gold clubs having a full length of 45 inches. Each of ten right-hitting golfers having a handicap of 15 to 30 struck 10 golf balls with each club. There were measured the flight distance and the amount of swerve to the right or left of the stopping position of a struck ball with respect to the target flight direction, and the standard deviation values thereof were calculated. The amount of swerve was shown by a positive value for the both cases of swerving to the right and the left. Each of the results of measurement of the flight distance and the amount of swerve shown in Table 1 is the average of found values obtained by striking 100 balls (10 balls×10 golfers) for each club. The larger the value, the better the flight distance performance. The smaller the value, the better the direction performance.
FIG. 7
The test results are shown in Table 1 and .
TABLE 1
Com.
Com.
Ex. 1
Ex. 1
Ex. 2
Ex. 3
Ex. 2
Ex. 4
Ex. 5
Moment of inertia Ix about vertical axis
5600
5400
5400
5400
5400
5000
5000
(g · cm<sup>2</sup>)
Radius of curvature Rx of face bulge
25
20
20
20
20
20
20
(inch)
Moment of inertia Iy about horizontal
1500
2300
3200
3500
3500
4500
3500
axis (g · cm<sup>2</sup>)
Radius of curvature Ry of face roll
8
10
14
18
20
14
14
(inch)
Ry/Rx ratio
0.32
0.50
0.70
0.90
1.00
0.70
0.70
Rx/Ix ratio
0.0045
0.0037
0.0037
0.0037
0.0037
0.0040
0.0040
Ry/Iy ratio
0.0053
0.0043
0.0044
0.0051
0.0057
0.0031
0.0040
Iy/Ix ratio
0.27
0.43
0.59
0.65
0.65
0.90
0.70
(Ry/Rx)/(Iy/Ix) ratio
1.2
1.2
1.2
1.4
1.5
0.8
1.0
Thickness of center thick wall region of
3.3
3.3
3.4
3.5
3.7
3.4
3.4
face portion (mm)
Flight distance (standard deviation A)
23.7
19.1
17.3
17.6
21.1
20.2
17.9
(yard)
Flight directionality (standard deviation
14.1
12.7
11.8
11.2
12.0
14.1
12.4
B of the amount of swerve)
(A + B)/2
37.8
31.8
29.1
28.8
33.1
34.3
30.3
Ex. 6
Ex. 7
Ex. 8
Ex. 9
Ex. 10
Ex. 11
Moment of inertia Ix about vertical axis
5000
5400
5400
5000
5900
4000
(g · cm<sup>2</sup>)
Radius of curvature Rx of face bulge
20
20
15
25
18
25
(inch)
Moment of inertia Iy about horizontal
1750
1600
3200
3000
3600
2000
axis (g · cm<sup>2</sup>)
Radius of curvature Ry of face roll
14
13
10
18
10
18
(inch)
Ry/Rx ratio
0.70
0.65
0.67
0.72
0.56
0.72
Rx/Ix ratio
0.0040
0.0037
0.0028
0.0050
0.0031
0.0063
Ry/Iy ratio
0.0080
0.0081
0.0031
0.0060
0.0028
0.0090
Iy/Ix ratio
0.35
0.30
0.59
0.60
0.61
0.50
(Ry/Rx)/(Iy/Ix) ratio
2.0
2.2
1.1
1.2
0.9
1.4
Thickness of center thick wall region of
3.4
3.4
3.2
3.6
3.3
3.6
face portion (mm)
Flight distance (standard deviation A)
18.8
19.7
20.1
19.2
20.5
21.7
(yard)
Flight directionality (standard deviation
13.5
14.6
15.4
15.3
14.9
15.0
B of the amount of swerve)
(A + B)/2
32.3
34.3
35.5
34.5
35.4
36.7
From the results shown in Table 1, it is confirmed that the golf club heads of the Examples according to the present invention have stable flight distance and flight direction performances as compared with the club heads of the Comparative Examples.
BRIEF DESCRIPTION OF THE DRAWINGS
FIG. 1
is a perspective view of a golf club head showing an embodiment of the present invention;
FIG. 2
FIG. 1
is a front view of the club head of ;
FIG. 3
FIG. 1
is a plan view of the club head of ;
FIG. 4
FIG. 3
is an enlarged cross sectional view along the line A-A of ;
FIG. 5
is a horizontal cross sectional view of a face portion of the club head at a horizontal plane passing through the sweet spot;
FIG. 6
FIG. 4
is a partially enlarged view of ;
FIG. 7
is a graph showing a relationship between the Iy/Ix ratio and the Ry/Rx ratio;
FIGS. 8A to 8C
are schematic plan views for illustrating the horizontal gear effect; and
FIGS. 9A and 9B
are cross sectional views for illustrating the vertical gear effect. | |
FIELD OF THE INVENTION
BACKGROUND ART
DESCRIPTION OF THE INVENTION
INDUSTRIAL APPLICABILITY
Example According to the Invention
The invention relates to a decorative element containing a gemstone and an electrically conductive layer on at least a partial area of the gemstone. The decorative element is suitable for electronic function control.
To date, gemstones have been employed almost exclusively for purely aesthetic purposes in accessories and on textiles, but hardly had any functional effect. In the field of wearable electronics (so-called “wearable technologies”), a market with enormous growth opportunities, they are lacking, because this field is associated by the users with functionality rather than decoration. The use of functional gemstones is generally conceivable wherever functionality and aesthetics are required, but especially in this case, the function control of electronic devices is a challenge. Touch-sensitive electronic sensors, such as those known from touchscreens, enable a comfortable function control of electronic devices using a finger or stylus. The input interface of an electronic device is the device portion by the touch of which a function is triggered. Gemstones serving as an input interface of devices and enabling an exact touch-sensitive handling of the devices are lacking.
The patent specification U.S. Pat. No. 7,932,893 describes a watch with touch-sensitive sensors serving to control a computer cursor.
The patent specification U.S. Pat. No. 6,868,046 discloses a watch with capacitive keys. The capacitive keys are operated manually using a finger and serve to control the hands of the watch.
WO2010/075599A1 describes a body made of a transparent material coated with a transparent electrically conductive layer. With the transparent electrically conductive layer, a contact to an inorganic semiconductor chip, an LED, is created.
US2006/0280040A1 discloses a watch in which a transparent electrically conductive layer is coated on the inside of the watch glass, and enamel is coated on the edge. This structure is supposed to avoid shortcuts.
EP 1 544 178 A1 describes a transparent element with transparent electrodes and a multi-layer structure. One of the layers is a transparent electrically conductive layer. The layers are coated on top of one another and also serve as an antireflection coating.
FR 1221561 discloses a decorative element that can be caused to light up by a phosphorescent material.
To date, there has been a lack of any technical solution of adapting gemstones for being suitable as a function control of electronic devices. It is the object of the present invention to provide a decorative element that enables function control of electronic devices.
A first subject matter of the present invention relates to a decorative element containing
(a) a gemstone,
(b) an electrically conductive layer on at least a partial area of the gemstone; and
(c) an electronic sensor.
The present invention further relates to a process for the function control of electronic devices, comprising the following steps:
(a) providing a gemstone coated with an electrically conductive layer;
(b) touching the electrically conductive layer with a finger or stylus; and
(c) triggering a function of the evaluation sensor system by the touch.
Gemstones coated with an electrically conductive layer can be used for the function control of electronic devices, for example, after being adhesively bonded to capacitive touchscreens, touchpads or trackpads, as are known, for example, from smartphones or laptops. Preferably, the electronic sensors for function control are contained in said touchscreens, touchpads or trackpads. In further possible applications preferred according to the invention, the electrically conductive layer is directly connected with the evaluation sensor system (see below). Decorative elements having such a structure are suitable for a wide range of applications. In jewels, for example, bracelets, rings, necklaces or brooches, that contain electronic components/devices, the decorative element according to the invention preferably comprises an additional electronic sensor for function control. The invention also relates to objects containing a decorative element according to the invention. For example, the decorative element may be advantageously incorporated in so-called “activity trackers”, to which the invention thus also relates. Further possible applications are mentioned in the following. The present invention also relates to the use of the decorative element according to the invention for the function control of electronic devices.
Surprisingly, it has been found that a combination of a gemstone with an electrically conductive layer and an electronic sensor as an input interface is suitable for a wide variety of purposes, for example, for controlling the brightness of a display. The combination according to the invention provides for a wide variety of possible applications in the design and technology fields, both as an input interface for function control and as a gemstone. A gemstone with an electrically conductive layer enables the function control of electronic devices by touching the electrically conductive layer with a finger or an electrically conductive stylus. Layers made of a wide variety of materials may be applied as the electrically conductive layer, for example, metallic layers or transparent conductive oxide layers are suitable (see below). Preferably, the electrically conductive layer is transparent in order that the color of the gemstone can be perceived. Preferably according to the invention, the gemstone is transparent in order to obtain brilliancy. From a brilliant appearance, gemstones obtain a particularly aesthetic effect. A combination of a transparent electrically conductive layer and a transparent gemstone is particularly preferred according to the invention.
Preferably, a wavelength-selective layer is additionally applied to at least a partial area of the gemstone. The combination of a transparent gemstone with a transparent electrically conductive layer and a wavelength-selective layer improves the brilliancy of the gemstone. The term “transparency” means the ability of matter to transmit electromagnetic waves (transmission). If a material is transparent for incident electromagnetic radiation (photons) of a more or less wide frequency range, the radiation can penetrate the material almost completely, i.e., it is hardly reflected and hardly absorbed. Preferably according to the invention, “transparency” means a transmission of at least 60% of the incident light, preferably more than 70%, more preferably more than 80%. A so-called mirror coating may also be applied to gemstones. A mirror coating is, for example, a silver layer applied by a wet chemical process. If appropriate, a mirror coating may also serve as a wavelength-selective layer.
Preferably according to the invention, facetted gemstones are employed. According to the invention, “faceting” means the design of a surface of a gemstone with polygons or so-called n-gons (n≥3); facets are usually obtained by grinding a rough crystal, but are also available by pressing methods.
Preferably according to the invention, the gemstones have a plano-convex or plano-convex-concave geometry, since the gemstones can be applied very readily because of the planar region. The curved convex or convex-concave region facilitates a comfortable operation of the decorative component, and the aesthetic of the decorative element is supported. The terms “convex” and “concave” relate to an imaginary enveloping area above or below the facets, and the definitions are to be understood by analogy with lenses in optics. The convex and concave regions may be either symmetrical or asymmetrical.
FIGS. 1
2
2
3
3
a
d
a
b
Possible structures of the decorative element (composite body) are shown in ), () to (), and () to (), the reference symbols having the following meanings:
1
() decorative element;
2
() gemstone;
3
() wavelength-selective coating;
4
() adhesive;
5
5
1
5
2
5
3
5
4
5
5
5
6
(.), (.), (.), (.), (.) and (.) are partial areas with an electrically conductive layer;
() electrically conductive layer;
6
() electrically conductive connection;
7
() evaluation sensor system;
8
() contacting by pogo pin;
9
() contacting by gemstone setting;
10
() contacting by electrically conductive adhesive;
11
() contacting by electrically conductive film, metallic mirroring, or electrically conductive elastomer;
12
() touch with a finger or stylus;
13
() movement in the direction of the arrow;
14
() movement in the altered direction of the arrow.
2
2
Preferably, the decorative element comprises a transparent faceted gemstone of glass and a transparent electrically conductive layer. Particularly preferred is an embodiment of the decorative element with a transparent faceted gemstone of glass with indium tin oxide as a transparent electrically conductive layer, and additionally a wavelength-selective layer composed of a sequence of SiOand TiOlayers. In a more particularly preferred embodiment of the decorative element, an electronic sensor is additionally contained.
FIG. 1
FIG. 1
FIG. 1
5
Preferably according to the invention, the electrically conductive layer (see below) is applied to the preferred convex/concave curved surface of the gemstone (). In , the layer () is drawn in discontinuity, because it may also be deposited in spatially separated regions according to the invention (see below). The wavelength-selective layer (see below) is preferably provided directly on the planar side opposite to the convex/concave curved surface, i.e., on the backside of the plano-convex-concave or plano-convex gemstone ().
Connecting the decorative element with an evaluation sensor system (see below) enables the function control of electronic devices. Touching the electrically conductive layer with a finger or an electrically conductive stylus triggers a signal that serves for the function control of electronic devices. Especially for wearable electronic devices, the function control of the electronic devices is a challenge because of their small size. The decorative element according to the invention combines a high brilliancy with the function of an input interface. However, the decorative element according to the invention is not limited to the field of wearable electronic devices. Its use is conceivable wherever functionality and aesthetics are desirable. For example, the use of the decorative element as a light switch is conceivable. In such a case, the light is switched on or off by touching the decorative element, and a dimming function is also possible. Generally, functions of electronic devices, such as television sets, radios or media players, can be controlled by the decorative element.
Another application of the decorative element is represented, for example, by rings and earrings, in which it serves as a gemstone and at the same time enables the function control of an electronic device. Thus, for example, a ring equipped with the decorative element according to the invention can be employed for the measurement of particular body functions. A wide variety of function control possibilities are conceivable, for example, a switching on and off function, or possible switching between different operation modes.
The decorative element according to the invention may also be employed for the function control for so-called switchable effects, for example, for a color change, for example, by means of RGB-LEDs, a gemstone or, for example, the display functions of a so-called smart watch. Switchable effects can be controlled with the decorative element and a suitable evaluation sensor system (see below) by touching the electrically conductive layer of the decorative element, for example, with a finger. Touching the electrically conductive layer of the decorative element may cause, for example, the color change of a gemstone.
The decorative element or a plurality of decorative elements may be integrated, for example, into a bracelet, in order to control the functions of, for example, a smart watch or an activity sensor (activity tracker). When a plurality of decorative elements is provided, the individual decorative elements can be employed by themselves for function control. The decorative elements can also be connected with each other, for example, by cables, for function control, so that only the successive touching of several gemstones causes a function (see below), for example, the brightness regulation of a display, and the volume regulation of speakers.
Gemstone
The gemstone can be made of a wide variety of materials, for example, glass, plastic, ceramic or gems or semi-precious stones. Gemstones made of glass or plastic are preferred according to the invention, because they are lowest cost and are most readily provided with facets. The use of glass is particularly preferred according to the invention. Faceted gemstones of glass are particularly preferred. The gemstones preferably comprise convex curved or convex-concave curved regions. This means that concave curved regions may also be present in addition to the convex curved regions on the faceted side. The side of the gemstone opposite the faceted side is either planar (preferably) or else concave. Particularly preferred are faceted gemstones of convex, especially plano-convex, geometry.
Glass
The invention is not limited in principle with respect to the composition of the glass. “Glass” means a frozen supercooled liquid that forms an amorphous solid. According to the invention, both oxidic glasses and chalcogenide glasses, metallic glasses or non-metallic glasses can be employed. Oxynitride glasses may also be suitable. The glasses may be one-component (e.g., quartz glass) or two-component (e.g., alkali borate glass) or multicomponent (soda lime glass) glasses. The glass can be prepared by melting, by sol-gel processes, or by shock waves. The methods are known to the skilled person. Inorganic glasses, especially oxidic glasses, are preferred according to the invention. These include silicate glasses, borate glasses or phosphate glasses. Lead-free glasses are particularly preferred.
2
For the preparation of the gemstones, silicate glasses are preferred. Silicate glasses have in common that their network is mainly formed by silicon dioxide (SiO). By adding further oxides, such as alumina or various alkali oxides, alumosilicate or alkali silicate glasses are formed. If phosphorus pentoxide or boron trioxide are the main network formers of a glass, it is referred to as a phosphate or borate glass, respectively, whose properties can also be adjusted by adding further oxides. These glasses can also be employed according to the invention. The mentioned glasses mainly consist of oxides, which is why they are generically referred to as oxidic glasses.
In a preferred embodiment according to the invention, the glass composition contains the following components:
2
(a) about 35 to about 85% by weight SiO;
2
(b) 0 to about 20% by weight KO;
2
(c) 0 to about 20% by weight NaO;
2
(d) 0 to about 5% by weight LiO;
(e) 0 to about 13% by weight ZnO;
(f) 0 to about 11% by weight CaO;
(g) 0 to about 7% by weight MgO;
(h) 0 to about 10% by weight BaO;
2
3
(i) 0 to about 4% by weight AlO;
2
(j) 0 to about 5% by weight ZrO;
2
3
(k) 0 to about 6% by weight BO;
(l) 0 to about 3% by weight F;
(m) 0 to about 2.5% by weight Cl.
All stated amounts are to be understood as giving a total sum of 100% by weight, optionally together with further components. The faceting of the gemstones is usually obtained by grinding and polishing techniques that are adequately familiar to the skilled person.
For example, a lead-free glass, especially the glass used by the company Swarovski for Chessboard Flat Backs (catalogue No. 2493), which shows a transmission of >95% in the range of 380-1200 nm, is suitable according to the invention.
Plastic
acrylic glass (polymethyl methacrylates, PMMA),
polycarbonate (PC),
polyvinyl chloride (PVC),
polystyrene (PS),
polyphenylene ether (PPO),
polyethylene (PE),
poly-N-methylmethacrylimide (PMMI).
As another raw material for the preparation of the gemstone (a), plastics can be employed. Transparent plastics are preferred according to the invention. Among others, the following materials are suitable according to the invention:
The advantages of the plastics over glass reside, in particular, in the lower specific weight, which is only about half that of glass. Other material properties may also be selectively adjusted. In addition, plastics are often more readily processed as compared to glass. Drawbacks include the low modulus of elasticity and the low surface hardness as well as the massive drop in strength at temperatures from about 70° C., as compared to glass. A preferred plastic according to the invention is poly-N-methylmethacrylimide, which is sold, for example, by Evonik under the name Pleximid® TT70. Pleximid® TT70 has a refractive index of 1.54, and a transmittance of 91% as measured according to ISO 13468-2 using D65 standard light.
Geometry
The geometric design of the gemstone is not limited in principle and predominantly depends on design aspects. The gemstone is preferably square, rectangular or round. Preferably according to the invention, the gemstone is faceted. The gemstone preferably has a convex, especially a plano-convex geometry. Preferably, the gemstone contains a plurality of facets on the preferably convex-curved side.
The type of faceting is closely related to the geometry of the gemstone. In principle, the geometric shape of the facets is not limited. Preferred according to the invention are square or rectangular facets, especially in combination with a gemstone with square or rectangular dimensions and a plano-convex geometry. However, faceted gemstones that are round may also be used.
Sensors
The function control of electronic devices using a finger or electrically conductive stylus is efficiently enabled by touch-sensitive electronic circuitry, as employed, for example, for touchscreens. Different electronic sensors are suitable for touch-sensitive electronic circuitry. Preferably according to the invention, resistive or capacitive sensors, more preferably capacitive sensors, are used as electronic sensors. Capacitive sensors include an electronic component with a capacitor and an input interface. In the decorative element, the input interface is the gemstone with the electrically conductive layer. Upon touching the input interface with a finger or an electrically conductive stylus, the capacitor changes its capacitance. This change is detected electronically and processed further by means of further electronic control elements. The capacitive or resistive sensors and the further processing electronic control elements are referred to as “evaluation sensor system”. Preferably, the electrically conductive layer is connected with the evaluation sensor system. This enables a very good operability of the decorative element.
The embodiments of capacitive and resistive sensors are adequately familiar to the skilled person. Resistive sensors comprise an electronic component with two separate electric contact surfaces. The two separate electric contact surfaces are connected by the gemstone coated with the electrically conductive layer. In the resistive sensors, touching the coated gemstone with a finger or with an electrically conductive stylus results in a current flow. The current flow can be detected electronically. The detection of the current flow as a control signal enables the function control of electronic devices. A Darlington transistor, for example, is suitable as a resistive sensor.
FIGS. 1 and 2
a
FIG. 2
b
FIG. 2
a
d
2
8
9
The connection between the input interface and the sensors is preferably created by an electrically conductive contact (to ). This has the advantage that the function control is not adversely affected. According to the invention, an electrically conductive contact is possible, for example, by using a pogo pin (). The pogo pin () creates an electrically conductive connection between the electrically conductive layer(s) on the gemstone and the sensor(s) by spring pressure onto the electrically conductive layer. Alternatively, an electrically conductive gemstone setting may also be used for the contacting (). For example, an electrically conductive portion of the gemstone setting () serves for holding the gemstone. The connection between the electrically conductive layer and the electrically conductive portion of the gemstone setting creates the contacting.
c
FIG. 2
d
FIG. 2
d
FIG. 2
Alternatively, an electrically conductive adhesive (), for example, 3M™ 5303 R-25μ/5303 R-50μ from the 3M company, an electrically conductive adhesive sheet (), for example, 3M® Anisotropic Conductive Film 7379 from the 3M company, or an electrically conductive elastomer (), for example, Silver Zebra® Connector from the company Fuji Polymer Industries Co. Ltd., for example, are suitable as the electrically conductive contacting. The electrically conductive connection may also be created by a wire connection. The possibilities of electrically conductive connection are adequately familiar to the skilled person.
Push-Type and Slide-Type Input
a
b
FIGS. 3and 3
The function control by means of an electrically conductive layer is possible in different ways. One embodiment is push-type input. In push-type input, a function of the evaluation sensor system, for example, the switching on or off of an electronic device, is triggered by the touch of the electrically conductive layer with a finger or an electrically conductive stylus. For push-type input, it is not required that the whole surface of the gemstone is coated with the electrically conductive layer. The electrically conductive layer may also be coated only on a partial area of the gemstone surface ().
a
b
FIGS. 3and 3
FIG. 1
For example, if the electrically conductive layer is applied to at least two electrically separated regions of the gemstone surface (dashed rectangles in ) and if the separated regions cause different functions, then electrically conductive contacting between the regions of the electrically conductive layer and the evaluation sensor system is required (). For example, one region enables the electronic device to be switched on and off, while the other region enables switching between the operational modes, for example. This results in a larger number of possibilities for function control. Since decorative elements are often incorporated in a setting, the connection to the electrically conductive contacting can be effected, for example, in the setting (see above). The different possibilities of preparing an electrically conductive connection are adequately familiar to the skilled person (see above).
5
1
5
2
5
3
5
4
5
5
5
6
12
13
14
a
FIG. 3
b
FIG. 3
a
b
FIGS. 3and 3
a
b
FIGS. 3and 3
Slide-type input is another possibility of function control. In this type of input, it is required that the electrically conductive layer is applied to at least two separated regions of the curved faceted surface (dashed rectangles . and . in as well as ., ., . and . in ). Function control is effected by a predefined succession of touches of the separated regions with a finger of with an electrically conductive stylus ( in ). The finger or electrically conductive stylus moves in the direction of the arrow ( and , respectively, in ). This comfortable type of input is also known from smartphones. The electrically conductive connection between the electrically conductive layer and the evaluation sensor system can be effected in different ways that are adequately familiar to the skilled person (see above).
5
3
5
4
5
6
5
5
b
FIG. 3
Therefore, in both push-type and slide-type input, the electrically conductive layer in at least two separated regions is of advantage for a comfortable function control. Therefore, the electrically conductive layer is preferably applied to at least two separated regions of the curved faceted gemstone surface. Further possibilities of function control are obtained if the push-type and slide-type input is combined in one decorative element, for example, slide-type input with the space-separated regions ., ., ., and push-type input with the region . ().
Decorative elements according to the invention that have push-type and/or slide-type input can be employed, for example, in bracelets, rings, necklaces, brooches, pockets, headsets or activity trackers. The jewels, such as bracelets, rings, necklaces or brooches, either contain themselves electronic devices with switchable functions, such as light control capabilities, or are used as a remote control for smartphones, headsets or activity trackers, for example. In a smartphone, for example, a function control is possible in which calls are accepted or rejected by touching the decorative element, volume regulation is conceivable for a headset, and switching between the operational modes in an activity tracker. The fields of application and the possibilities of function control are mentioned merely in an exemplary way, while a wide variety of controllable functions can be realized.
Electrically Conductive Layer
FIG. 1
In connection with an evaluation sensor system, the electrically conductive layer enables the function control of electronic devices. Preferably according to the invention, it is applied to the curved surface of the gemstone (), in order to enable a simple touch with a finger or with an electrically conductive stylus. The transmission properties of the electrically conductive layer affect the brilliancy of the decorative element. Therefore, the electrically conductive layer is preferably transparent within a range of 380 to 780 nm.
Because of their electrical conductivity, metallic layers are suitable as the electrically conductive layer. They can be deposited on the gemstone by suitable coating methods, for example, sputtering (see below). Metals like Cr, Ti, Zr, V, Mo, Ta and W are suitable for this. Metals like Al, Cu or Ag are less advantageous as the electrically conductive layer because of their lower chemical stability, but still suitable in principle. Chemical compounds with electrical conductivity properties may also be used as the electrically conductive layer, particularly chemical nitride compounds, for example, TiN, TiAlN or CrN. The metallic layers and electrically conductive chemical compounds are adequately familiar to the skilled person.
Transparent electrically conductive oxide layers can also be employed as the electrically conductive layer. They are well known to the skilled person. Transparent electrically conductive oxide layers have a good mechanical abrasion resistance, a good chemical resistance, and a good thermal stability. They contain semiconductive oxides. The semiconductive oxides obtain metallic conductivity from a suitable n doping. The transparent electrically conductive oxide layers are important components for transparent electrodes, for example, in flat screens or thin layer solar cells.
2
3
2
Indium tin oxide is the transparent electrically conductive oxide layer that is most readily technically accessible. It is a commercially available mixed oxide of about 90% InOand about 10% SnO. Indium tin oxide has very good transmission properties, a very good mechanical abrasion resistance, and a very good chemical resistance. Preferably according to the invention, indium tin oxide is used as the electrically conductive layer, and in particular, indium tin oxide is applied at a layer thickness of at least 4 nm to obtain electrical conductivity.
Aluminum-doped zinc oxide as the transparent electrically conductive oxide layer has good transmission properties and a good mechanical abrasion resistance. It is employed on an industrial scale, for example, in the field of solar technology. Further suitable transparent electrically conductive oxide layers include doped zinc oxides, such as gallium zinc oxide or titanium zinc oxide, doped tin oxides, such as fluorine-doped tin oxide, antimony tin oxide, or tantalum tin oxide, or doped titanium niobium oxide.
Preferably according to the invention, the electrically conductive layer comprises at least one component selected from the group of Cr, Ti, Zr, indium tin oxide, aluminum-doped zinc oxide, gallium zinc oxide, titanium zinc oxide, fluorine-doped tin oxide, antimony tin oxide, tantalum tin oxide, or titanium niobium oxide, or any combination of these components in any sequence of layers. More preferably, indium tin oxide is deposited for the decorative element according to the invention.
The methods for preparing electrically conductive layers are adequately familiar to the skilled person. These include, without limitation, PVD (physical vapor deposition) and CVD (chemical vapor deposition) methods. PVD methods are preferred according to the invention.
The PVD methods are a group of vacuum-based coating methods or thin layer technologies that are adequately familiar to the skilled person, being employed, in particular, for the coating of glass and plastic in the optical and jewelry industries. In the PVD process, the coating material is transferred into the gas phase. The gaseous material is subsequently led to the substrate to be coated, where it condenses and forms the target layer. With some of these PVD methods (magnetron sputtering, laser beam evaporation, thermal vapor deposition, etc.), very low process temperatures can be realized. In this way, a large number of coating materials can be deposited in a very pure form in thin layers. If the process is performed in the presence of reactive gases, such as oxygen, metal oxides, for example, may also be deposited. A preferred method according to the invention is a coating process by sputtering, for example, with the device Radiance from the company Evatec. Depending on the requirements of function an optical appearance, a typical layer system can consist of only one layer, but also of a large number of layers.
For the preparation of the separated regions of the electrically conductive layer on the curved faceted surface (see above), the gemstone is covered by a mask. The mask leaves the regions of the curved faceted surface exposed, on which the electrically conductive layer is deposited. Covers of plastic or metal are suitable as the mask, for example. An alternative possibility for preparing the separated regions of the electrically conductive layer on the curved faceted surface is cutting through this layer by means of a laser, for example, an Nd:YAG laser or an ultrashort pulse laser. The use of a laser enables a very precise preparation of the separated regions. The separation of the electrically conductive layer may also be effected by etching. Etching includes the application of a mask to the electrically conductive layer, for example, by using a photoresist. The etching creates the desired spatially separated regions of the electrically conductive layer. The photoresist is subsequently removed, for example, by wet chemical methods. The methods are adequately familiar to the skilled person.
Wavelength-Selective Layer
The wavelength-selective layer increases the brilliancy of the decorative element. The optional wavelength-selective layer is preferably provided between the gemstone and the evaluation sensor system. Preferably according to the invention, it will be realized in two different ways: by a wavelength-selective film or a wavelength-selective coating, which is prepared by PVD, CVD or wet-chemical methods. However, a wavelength-selective layer may also be obtained from a microstructured surface. The methods of microstructuring are well known to the skilled person.
As a result of the reflection of a defined range (=filtering) of the visible spectrum, the optical element gains brilliance and appears in a particular color to the viewer. The brilliance is additionally supported by the faceting of the gemstone. In a preferred embodiment of the invention, the wavelength-selective layer reflects a fraction of the light in the range of 380 to 780 nm, i.e., predominantly in the visible range.
a
b
FIGS. 5and 5
The wavelength-selective layer shows angle-dependent reflection (). The reflection interval is shifted as a function of the angle of incidence of the light onto the decorative element. Depending on the position of the decorative element, different color fractions are reflected.
In order to enable bonding of the individual components of the decorative element with UV-curing adhesives, the wavelength-selective layer is preferably at least partially transparent to UV light.
Preferably according to the invention, the wavelength-selective layer is a dielectric in order to enable unrestricted function control for separated regions of the electrically conductive layer (see above). If the wavelength-selective layer is electrically conductive, fault currents may occur.
According to the invention, the wavelength-selective layer could be applied on the gemstone surface between the electrically conductive layer and the gemstone surface in principle; however, this is one of the less preferred embodiments because of a possible reduction of brilliancy. If the wavelength-selective layer is applied on the planar side of the gemstone, there are multiple reflections within the gemstone, which lead to an increase of brilliancy.
Wavelength-Selective Films
Wavelength-selective films are commercially available under the designation “Radiant Light Film”. These are multilayered polymeric films that can be applied to other materials. These optical films are Bragg mirrors and reflect a high proportion of the visible light and produce brilliant color effects. A relief-like microstructure within a range of several hundred nanometers reflects the different wavelengths of the light, and interference phenomena occur, the colors changing as a function of the viewing angle.
Particularly preferred films according to the invention consist of multilayered polymeric films whose outermost layer is a polyester. Such films are sold, for example, by the company 3M under the name Radiant Color Film CM 500 and CM 590. The films have a reflection interval of 590-740 nm or 500-700 nm.
The wavelength-selective film is preferably bonded with the gemstone by means of an adhesive. When the electrically conductive layer and the gemstone are transparent, the adhesive should also be transparent. In a preferred embodiment, the refractive index of the adhesive deviates by less than ±20% from the refractive index of the transparent gemstone. In a particular preferred embodiment, the deviation is <10%, even more preferably <5%. This is the only way to ensure that reflection losses because of the different refractive indices can be minimized. The refractive indices can also be matched to one another by roughening the respective boundary layers (moth eye effect). So-called “moth eye surfaces” consist of fine nap structures that change the refraction behavior of the light, not suddenly, but continuously in the ideal case. The sharp boundaries between the different refractive indices are removed thereby, so that the transition is almost fluent, and the light can pass through unhindered. The structural sizes required for this must be smaller than 300 nm. Moth eye effects ensure that the reflection at the boundary layers is minimized, and thus a higher light yield is achieved in the passage through the boundary layers.
Adhesives that can be cured by means of UV radiation are preferred according to the invention. Both the UV-curing adhesives and the methods for determining the refractive index are well known to the skilled person. Particularly preferred according to the invention is the use of acrylate adhesives, especially of modified urethane acrylate adhesives. These are sold by numerous companies, for example, by Delo under the designation Delo-Photobond® PB 437, an adhesive that can be cured by UV light within a range of 320-42 nm.
Wavelength-Selective Coating
The coating materials are well known to the skilled person. In a preferred embodiment of the invention, the wavelength-selective coating is a dielectric (see above). Dielectric coating materials can be applied to the gemstone by one of the common coating methods. Successive layers of different dielectric materials can also be applied. The methods of preparing coatings and the coatings themselves are adequately known to the skilled person. These include, among others, PVD (physical vapor deposition) methods, CVD (chemical vapor deposition) methods, paint-coating methods and wet chemical methods according to the prior art. PVD methods are preferred according to the invention (see above).
2
2
3
2
3
3
2
3
4
2
5
2
2
2
For the construction of a dielectric wavelength-selective coating according to the invention, the following coating materials are preferably suitable: MgF, SiO, CeF, AlO, CeO, ZrO, SiN, TaO, TiO, or any combination of such compounds in any sequence of layers, a succession of TiOand SiOlayers being particularly preferred. The desired degree of reflection and transmission can be adjusted by appropriately selecting the coating materials, number of layers and the layer thicknesses.
For the PVD layer production, a wide variety of commercial machines are available, for example, model BAK1101 from the company Evatec.
The decorative elements according to the invention can be employed for the function control of numerous electronic devices.
An example according to the invention with a gemstone and an electrically conductive layer was prepared.
Gemstone:
The non-mirrored Chessboard Flat Back 2493 (30 mm×30 mm) of the company D. Swarovski KG was used as a gemstone of glass.
Geometry:
The gemstone was a faceted solid with 30 mm edge length and a square base area with slightly rounded corners. The faceted upper part included convex curved areas. The total height of the solid was about 8 mm, the corner edge height was about 2.7 mm.
Transparent Electrically Conductive Layer:
The gemstone was coated with the transparent electrically conductive layer of indium tin oxide on the entire surface. The coating process was performed by sputtering with the PVD plant FHRline 400 of the company FHR.
2
3
2
2
−3
In order to improve the electrical and chemical properties and the mechanical abrasion resistance, the gemstone was first treated by ion etching in the plant FHRline 400. Thereafter, the sample was heated at a temperature of about 550° C. for about 30 minutes in the same plant FHRline 400. This was followed by the coating of the optical element with indium tin oxide in the same plant FHRline 400, wherein the mixed oxide had a customary ratio of about 90% InOto about 10% SnO. The pressure was about 3.3·10mbar, and the discharge power was about 1 kW. The layer thickness varied as a function of the surface geometry from about 140 nm to about 190 nm. The coating process was effected with using a protective gas of argon and 5 sccm O. Subsequently, the coated optical element was heated at a temperature of about 550° C. for about 20 minutes in the same plant FHRline 400.
Evaluation Sensor System and Structure of the Decorative Element:
The coated gemstone was connected on the backside by the pogo pin S7121-42R from the company Harwin Plc Europe with the circuit board Kingboard KB-6160 FR-4Y KB 1.55. The pogo pin S7121-42R was soldered with the circuit board. The distance between the circuit board and the coated backside of the gemstone was about 1.5 mm. The touch controller IQS228AS from the company Azoteq (Pty) Ltd. was used for function control. The touch controller IQS228AS was provided on the upper side of the circuit board between the gemstone and circuit board and was soldered with the circuit board. The touch controller IQS228AS was electrically connected with the pogo pin through a conducting path. The touch controller was supplied with power, and the signal of the touch controller transmitted, through a multi-pole cable. The structure was surrounded by a housing of polycarbonate of the type Makrolon® 2405 polycarbonate. The gemstone was connected at a distance of about 1.5 mm from the circuit board with the housing through the commercially available two-part epoxy resin adhesive 9030 CG 500 (A+B) 50 ml EUROPE/AMERICA, Material No. 5284198 from the company Swarovski. The housing had an inward running web of about 1.7 mm in order to enable a distance of the backside of the gemstone of about 1.5 mm from the circuit board, and the connection between the housing and gemstone by means of an adhesive. The multi-pole cable was led out of the housing through an opening in the housing.
In the following, the invention will be illustrated further by means of Examples and Figures without being limited thereto. The Figures show the following objects:
FIG. 1
: Structure of a decorative element. Electrically conductive layer in partial areas of the gemstone, and wavelength-selective coating on the planar side opposite the faceting.
a
FIG. 2
: Electrically conductive connection between the electrically conductive layer and the evaluation sensor system by a pogo pin.
b
FIG. 2
: Electrically conductive connection between the electrically conductive layer and the evaluation sensor system by an electrically conductive gemstone setting.
c
FIG. 2
: Electrically conductive connection between the electrically conductive layer and the evaluation sensor system by an electrically conductive adhesive.
d
FIG. 2
: Electrically conductive connection between the electrically conductive layer and the evaluation sensor system by an electrically conductive film or electrically conductive elastomer.
a
FIG. 3
: Decorative element with two separated regions of the transparent electrically conductive layer for push-type or slide-type input.
b
FIG. 3
: Decorative element with four separated regions of the transparent electrically conductive layer for push-type or slide-type input. | |
The debate opens amid ebbing political enthusiasm for banking union – originally planned as a three-stage process involving ECB bank supervision, alongside an agency to shut failing banks and a system of deposit guarantees. It would be the boldest step in European integration since the crisis. “We have to find a solution now,” said Michel Barnier, the EU Commissioner in charge of financial regulation, urging faster progress in the slow talks. “The next financial crisis is not going to wait for us.” ANGLO-GERMAN AXIS? In one sign of the divisions, Britain has repeatedly refused to sign off on the first pillar of the banking union framework, allowing the ECB monitor banks. Having earlier agreed, London now wants additional assurances from ministers this week that Britain, which is outside the euro and polices its own banks, will not face interference from the ECB-led euro bloc. Britain is likely to find a sympathetic ear in Berlin, which wants to keep London on side in its push to prevent stricter EU emissions rules to protect its luxury car makers. Before the ECB takes over as supervisor late next year, it will conduct health checks of the roughly 130 banks under its watch. This is the nub of the problem facing finance ministers at the two-day talks. With the euro zone barely out of recession, a failure to put aside money to deal with the problems revealed could rattle fragile investor confidence and compound borrowing difficulties for companies, potentially killing off the meek recovery. In turn, that raises the question about who pays for the holes that are found in balance sheets in countries such as Spain and Italy.
Europe stocks waver amid U.S. deadlock, China data
government shutdown and three days before the country is expected to reach its borrowing limit, unless lawmakers break a stalemate and raise the nations debt ceiling. On Sunday, Senate Republicans and Democrats leaders continued attempts to find a way to break the fiscal impasse between the Republican-led House and President Barack Obama. Read: Fed shutdown and your retirement: Remain calm Treasury Secretary Jack Lew has warned the U.S. will run out of borrowing authority on Oct. 17 unless Congress agrees to raise the debt ceiling. A failure to increase the limit could lead to a technical default, which some fear will drag the economy back into recession . U.S. stocks traded lower on Wall Street . China and Europe data Meanwhile in China, data over the weekend showed a surprise decline in exports in September , signaling the global economy is still struggling to recover. Additionally, Chinese consumer prices rose faster than expected in September, though remaining within the governments target range. On the data front in Europe, Eurostat said industrial production rebounded in August in the euro zone, rising 1% month-on-month. Industrial production data for August support our view that euro-area growth is resuming but is still weak, and risks remain skewed to the downside, said Fabio Fois, southern European economist at Barclays, in a note. Our tentative forecast for euro-area industrial production in Q3 (…) also points to a slowdown in economic activity. That said, we continue to expect euro-area GDP to have increased 0.2% q/q in Q3 (0.1pp below Q2), a view that is supported by various confidence data, including PMIs, he added. Among country-specific indexes, Frances CAC 40 index /quotes/zigman/3173214 FR:PX1 +0.07% inched 0.1% higher to 4,222.96, while Germanys DAX 30 index /quotes/zigman/2380246 DX:DAX -0.01% closed slightly lower at 8,723.81. | http://www.delcamp.us/rpt-europe-prepares-to-come-clean-on-hidden-bank-losses.html |
Conventionally, as a method of forming a thin film on a dielectric, silicon, or other substrate, for example, there is a thin film forming method by making a material of the thin film (hereafter, this is called a thin film material) into particulates, and making the particulate thin film material on the substrate. As the thin film forming methods, depending on difference between methods of making the thin film material into particulates, methods of making it deposited, and the like, for example, there are a sputtering method, a CVD (Chemical Vapor Deposition) method, an MBE (Molecular Beam Epitaxy) method, a laser ablation method, a vacuum deposition method, etc.
In addition, generally, thin film forming systems which are used when forming a thin film using the sputtering method etc. are classified into systems called a parallel plate type and an opposed type according to positional relation between a substrate, on which the thin film is formed, and a target for generating the particulate thin film material.
The parallel plate type thin film forming apparatus is arranged, for example, as shown in FIG. 13, so that a first principal surface 1A of the substrate 1, and the target 2B may be in parallel. At this time, the target 2B is mounted on a cathode 3, and by supplying electric power to the cathode 3, the particulate thin film material 2A is sputtered out of the target 2B between the target 2B and the substrate 1. Then, for example, by applying an electric field between the target 2B and the substrate 1, introducing the particulate thin film material 2A in a direction of the substrate 1 with accelerating it, and depositing the particulate thin film material 2A on the first principal surface 1A of the substrate 1, the thin film 2 is formed. In addition, at this time, the substrate 1 is fixed to a heat stage 13, and is heated from a backside 1B of the first principal surface 1A of the substrate 1 (hereafter, this is called a second principal surface) as shown in FIG. 13.
In the case of the parallel plate type thin film forming apparatus, the particulate thin film material 2A which is accelerated and is in a high energy state collides at a nearly vertical angle to the thin film formation surface of the substrate. Therefore, while the film formation speed of the thin film 2 is fast and productive efficiency is high, there is a problem that a damage to the surface of the thin film 2 deposited on the substrate 1 is serious. In order to make the damage to the surface of the thin film 2 small, for example, there is a method of making an acceleration of the particulate thin film material 2A small. Nevertheless, when an acceleration of the particulate thin film material 2A is made small, film formation speed of the thin film 2 drops, and productive efficiency drops. Then, in recent years, for example, the opposed type thin film forming apparatus is proposed as a thin film forming apparatus which replaces the parallel plate type thin film forming apparatus.
The opposed type thin film forming apparatus is arranged, for example, as shown in FIG. 14, so that two targets 2B may face each other on an extension in a direction parallel to an in-plane direction of the first principal surface 1A of the substrate 1. Also at this time, the each target 2B is mounted on a cathode 3, and by supplying electric power to the cathode 3, the particulate thin film material 2A is sputtered out of the targets 2B. Since the particulate thin film material 2A which is sputtered out of the each target 2B gather between the two targets 2B facing each other, when the particulate thin film material 2A is accelerated by applying an electric field between the two targets 2B, and are introduced on the first principal surface 1A of the substrate 1, the thin film material 2 is formed by the particulate thin film material 2A being deposited on the first principal surface 1A of the substrate 1.
In the case of the opposed type thin film forming apparatus, since an incident angle at the time of the particulate thin film material 2A colliding with the first principal surface 1A of the substrate 1 is as small as about 0° to 45°, damage which the thin film 2 deposited on the first principal surface 1A of the substrate 1 receives is small when the particulate thin film material 2A collides with the first principal surface 1A of the substrate 1. Therefore, since it is possible to introduce and deposit the particulate thin film material 2A on to the first principal surface 1A of the substrate 1 in a high energy state, and it is possible to form the thin film 2 damage of whose surface is small without reducing productive efficiency.
In addition, since other methods such as the CVD method, MBE method, laser ablation method, and other film formation methods also form a thin film with principles and apparatuses similar to the sputtering method, detailed description is omitted.
The parallel plate type and opposed type thin film forming systems are used, for example, when producing microwave devices such as a GPS (Global Positioning Systems) array antenna and a microwave integrated circuit. In the microwave device, for example, as shown in FIGS. 15 and 16, a circuit pattern 2C is provided on the first principal surface 1A of the substrate 1, and a ground plane 2D is provided on the second principal surface 1B of the substrate 1. Here, FIG. 16 is a sectional view taken on line D-D′ in FIG. 15.
The microwave device is operated using a change of a magnetic field generated in connection with a leakage electric field generated between the circuit pattern 2C and the ground plane 2D, for example, as shown in FIG. 17. At this time, when the circuit pattern 2C and the ground plane 2D are oxide layer superconductors, for example, it is possible to obtain smaller surface resistance and higher operating characteristics in comparison with usual conductors. Therefore, recently, various microwave devices using the oxide superconductors have attracted attention (for example, refer to S. 0hshima, “High-temperature superconducting passive microwave devices, filters, and antennas”, Supercond. Sci. Technol., 13, 2000, p. 103-108).
In a microwave device using the oxide superconductors, for example, a dielectric substrate such as magnesium oxide (MgO) or sapphire (A1203) is used for the substrate 1, and oxide superconductors such as YBCO or BSCCO are used for the circuit pattern 2C and the ground plane 2D.
When producing a microwave device using the oxide superconductors, first, as shown in FIG. 18, thin films 2C′ and 2D of the oxide superconductors are formed on the first principal surface 1A and the second principal surface 1B of the dielectric substrate 1, respectively. The parallel plate type and opposed type thin film forming apparatuses are used for formation of the thin films 2C′ and 2D. At this time, it is assumed that the target 2B is constructed of, for example, a material of YBa2CU3Ox, Y203, BaO, CuO, or the like which is used for formation of YBCO which is one kind of oxide superconductors. In addition, the dielectric substrate 1 is heated at, for example, about 800° C. at this time.
In addition, when forming the thin films 2C′ and 2D, for example, after forming the thin film 2C′ on the first principal surface 1A of the dielectric substrate 1, the dielectric substrate 1 is turned over, and the thin film 2D on the second principal surface 1B of the dielectric substrate 1 is formed. At this time, respective thin films 2C′ and 2D on the first principal surface 1A and second principal surface 1B of the dielectric substrate 1 are formed with composition of target 2B and conditions in an apparatus at the time of formation being fixed, for example, so as to become the same film quality and film thickness.
Next, as shown in FIG. 19, an etching resist 12 matched with the circuit pattern 2C is formed on one thin film, for example, the thin film 2C′ on the first principal surface 1A of the dielectric substrate 1. At this time, although illustration is omitted, a resist is formed, for example, also on the backside of the surface on which the etching resist 12 is formed, that is, the thin film 2D on the second principal surface 1B of the substrate 1. Then, unnecessary portions are removed by etching the thin film 2C′ on the surface on which the etching resist 12 is formed, and the circuit pattern 2C as shown in FIG. 15 is formed.
Nevertheless, when forming the thin films 2 in both sides of the substrate 1 by the conventional art, it is necessary to form single sides separately. Therefore, for example, while turning the substrate 1 over and forming the thin film 2 on the second principal surface 1B of the substrate 1 after forming the thin film 2 on the first principal surface 1A of the substrate 1, film quality of the thin film 2 formed on the first principal surface 1A of the substrate 1 may change. In particular, when forming the thin films 2C′ and 2D of the oxide superconductors like the microwave device, there was a problem that degradation of the film quality due to a timing change occurred easily.
In addition, even if thin films are formed under the same conditions using the same thin film forming apparatus, it is apt to generate difference between the film qualities of the thin film 2 formed at a time and the thin film 2 formed at a second time because of states of the target 2B, temperature unevenness at the time of heating, etc. Therefore, the conventional methods for forming a thin film had a problem that it was difficult to equalize the film qualities of the thin film of the first principal surface 1A and the thin film of the second principal surface 1B of the substrate 1.
In particular, the oxide superconductor used when producing the device is deficient in chemical stability. Therefore, when forming single sides separately at the time of forming the thin films 2C′ and 2D of the oxide superconductor on both sides of the substrate 1, degradation of the film quality and decrease of uniformity of the thin film 2C′ formed on the first principal surface 1A and the thin film 2D formed on the second principal surface 1B of the dielectric substrate 1 are apt to be generated. Therefore, for example, there was a problem that difference between electrical characteristics of the circuit pattern 2C, and electrical characteristics of the ground plane 2D arose and operation of the device became unstable.
In addition, upsizing of the substrate 1 used for manufacturing the microwave device etc. has been advancing recently. Therefore, when forming single sides separately at the time of forming thin films on both sides of the first principal surface 1A and the second principal surface 1B of the substrate 1, degradation and unevenness of film quality become remarkable. Furthermore, for example, there was a problem that time and energy consumption required for formation of the thin films increased.
Moreover, when the time required for the formation of the thin films became long, there was a problem that productive efficiency of the thin films dropped and manufacturing cost rose.
Hence, the present invention aims at providing a thin film forming method and a thin film forming apparatus which can reduce degradation and dispersion of film qualities of thin films of respective surfaces of the substrate when depositing a material, which is made into particulates, on both sides of a substrate, for example, when forming thin films of oxide superconductors or the like.
In addition, the present invention aims at providing a thin film forming method and a thin film forming apparatus which can reduce production cost at the time of forming thin films, such as oxides superconductors, on both sides of a substrate.
| |
PROBLEM TO BE SOLVED: To provide a vehicle cooling device which suppresses excessive cooling of an air suction gas containing an EGR gas, that circulates in an intercooler in an engine with supercharger and intercooler, and can suppress generation of condensed water.
SOLUTION: The vehicle cooling device comprises: an engine 4 having an EGR device 14, a supercharger 10, and an intercooler 12 for cooling an air suction gas; a cooling water flow channel 32 in which cooling water for cooling motor drive devices 6, 16, 18 flows; and a cooling water flow rate regulator 48. The cooling water flow channel 32 has: a first flow channel 32a in which first cooling water for cooling high voltage components 16, 18 of the motor drive device flows; and a second flow channel 32b, which branches from the first flow channel and in which second cooling water for cooling the intercooler flows. The cooling water flow rate regulator is provided in the second flow channel of the cooling water flow channel, and can regulate a flow rate of the second cooling water.
SELECTED DRAWING: Figure 2
COPYRIGHT: (C)2022,JPO&INPIT | |
Cotton String Balls (3 Pack)
This pack of Cotton String Balls contains three balls of string that weigh approximately 50g. Each ball of string is constructed from recycled white cotton and is suitable for a wide variety of jobs and tasks.
Use this popup to embed a mailing list sign up form. Alternatively use it as a simple call to action with a link to a product or a page.
By clicking enter you are verifying that you are old enough to consume alcohol. | https://makeitmagic.co.uk/products/cotton-string-balls-3-pack |
Number Theory in Function Fields [electronic resource] / by Michael Rosen
- Author:
- Rosen, Michael
- Published:
- New York, NY : Springer New York : Imprint: Springer, 2002.
- Physical Description:
- XI, 358 pages : online resource
- Additional Creators:
- SpringerLink (Online service)
Access Online
- Series:
- Contents:
- 1 Polynomials over Finite Fields -- 2 Primes, Arithmetic Functions, and the Zeta Function -- 3 The Reciprocity Law -- 4 Dirichlet L-Series and Primes in an Arithmetic Progression -- 5 Algebraic Function Fields and Global Function Fields -- 6 Weil Differentials and the Canonical Class -- 7 Extensions of Function Fields, Riemann-Hurwitz, and the ABC Theorem -- 8 Constant Field Extensions -- 9 Galois Extensions — Hecke and Artin L-Series -- 10 Artin’s Primitive Root Conjecture -- 11 The Behavior of the Class Group in Constant Field Extensions -- 12 Cyclotomic Function Fields -- 13 Drinfeld Modules: An Introduction -- 14 S-Units, S-Class Group, and the Corresponding L-Functions -- 15 The Brumer-Stark Conjecture -- 16 The Class Number Formulas in Quadratic and Cyclotomic Function Fields -- 17 Average Value Theorems in Function Fields -- Appendix: A Proof of the Function Field Riemann Hypothesis -- Author Index.
- Summary:
- Elementary number theory is concerned with arithmetic properties of the ring of integers. Early in the development of number theory, it was noticed that the ring of integers has many properties in common with the ring of polynomials over a finite field. The first part of this book illustrates this relationship by presenting, for example, analogues of the theorems of Fermat and Euler, Wilsons theorem, quadratic (and higher) reciprocity, the prime number theorem, and Dirichlets theorem on primes in an arithmetic progression. After presenting the required foundational material on function fields, the later chapters explore the analogy between global function fields and algebraic number fields. A variety of topics are presented, including: the ABC-conjecture, Artins conjecture on primitive roots, the Brumer-Stark conjecture, Drinfeld modules, class number formulae, and average value theorems. The first few chapters of this book are accessible to advanced undergraduates. The later chapters are designed for graduate students and professionals in mathematics and related fields who want to learn more about the very fruitful relationship between number theory in algebraic number fields and algebraic function fields. In this book many paths are set forth for future learning and exploration. Michael Rosen is Professor of Mathematics at Brown University, where hes been since 1962. He has published over 40 research papers and he is the co-author of A Classical Introduction to Modern Number Theory, with Kenneth Ireland. He received the Chauvenet Prize of the Mathematical Association of America in 1999 and the Philip J. Bray Teaching Award in 2001.
- Subject(s):
- ISBN:
- 9781475760460
- Digital File Characteristics:
- text file PDF
- Note:
- AVAILABLE ONLINE TO AUTHORIZED PSU USERS. | https://catalog.libraries.psu.edu/catalog/15203702 |
This application claims priority to Korean Patent Application No. 10-2013-0052617 filed on May 9, 2013, and all the benefits accruing therefrom under 35 U.S.C. §119, the disclosure of which is incorporated herein by reference in its entirety.
1. Field
The invention relates to a mask and a method of manufacturing the same.
2. Description of the Related Art
Flat panel display devices are replacing cathode-ray tube display devices, due to their lightweight and thin characteristics. Typical examples of the flat panel display devices include liquid crystal display devices (“LCDs”) and organic light-emitting diode display devices (“OLEDs”). OLEDs have excellent brightness and viewing angle characteristics as compared to LCDs, and can be realized as ultra-thin display devices because they do not require a separate light source such as a backlight unit.
OLEDs use a phenomenon in which electrons and holes injected into an organic thin film from a cathode and an anode are recombined to form excitons, and light having a specific wavelength is emitted by energy generated from the excitons.
OLEDs are classified into a passive-matrix type and an active-matrix type according to a driving method. Active-matrix OLEDs include a circuit using a thin-film transistor (“TFT”).
Passive-matrix OLEDs are easy to manufacture because their display region is constructed in a simple matrix of anodes and cathodes. However, the passive-matrix OLEDs are restricted in application fields of low-resolution and small displays due to problems with resolution, an increase in driving voltage and a decrease in material duration. Active-matrix OLEDs can provide stable luminance due to a constant current supplied to each pixel using a TFT located at each pixel of a display region. With their low power consumption, the active-matrix OLEDs can be implemented as high-resolution and relatively large-sized displays.
To realize a full-color OLED, red, green and blue light-emitting layers may be formed using a laser induced thermal imaging (“LITI”) method among various methods. In the LITI method, a laser beam emitted from a laser source is patterned using a mask having patterns, and the patterned laser beam is irradiated onto a donor substrate which includes a base substrate, a light-to-heat conversion layer and a transfer layer (e.g., an organic layer including a light-emitting layer), such that part of the transfer layer can be transferred onto a device substrate. Accordingly, organic film layer patterns including the light-emitting layer are formed on the device substrate. The LITI method is advantageous in that it can finely pattern each light-emitting layer and is a dry-etching process.
The mask having the patterns includes transmissive regions which transmit laser light and non-transmissive regions which do not transmit laser light.
One or more exemplary embodiment of the invention provides a method of manufacturing a mask in a simplified process and with improved process efficiency and a mask manufactured by the method.
One or more exemplary embodiment of the invention also provides a method of manufacturing a large-area mask more easily and a mask manufactured by the method.
However, embodiments of the invention are not restricted to the one set forth herein. The above and other features of the invention will become more apparent to one of ordinary skill in the art to which the invention pertains by referencing the detailed description of the invention given below. According to an exemplary embodiment of the invention, there is provided a method of manufacturing a mask. The method includes: providing a base substrate including light-absorbing layer patterns on a first surface thereof; providing a reflective layer on the light-absorbing layer patterns and the first surface of the base substrate; and providing reflective patterns by partially removing the reflective layer. The providing the reflective patterns includes removing the light-absorbing layer patterns and a portion of the reflective layer, by irradiating the light-absorbing layer patterns with laser light.
According to another exemplary embodiment of the invention, there is provided a method of manufacturing a mask. The method includes: providing a base substrate including light-absorbing layer patterns on a first surface thereof and a reflection preventing layer on a second surface opposite to the first surface thereof; providing a reflective layer on the light-absorbing layer patterns and the first surface of the base substrate; and providing reflective patterns by partially removing the reflective layer. The providing the reflective patterns includes removing the light-absorbing layer patterns and part of the reflective layer, by irradiating the light-absorbing layer patterns with laser light.
According to another exemplary embodiment of the invention, there is provided a mask including: reflective patterns on a first surface of a base substrate; and a reflection preventing layer covering a second surface opposite to the first surface of the base substrate. A side surface of the reflective patterns is uneven.
Hereinafter, exemplary embodiments of the invention will be described in detail with reference to the accompanying drawings. The features of the invention and methods for achieving the features will be apparent by referring to the exemplary embodiments to be described in detail with reference to the accompanying drawings. However, the invention is not limited to the exemplary embodiments disclosed hereinafter, but can be implemented in diverse forms. The matters defined in the description, such as the detailed construction and elements, are nothing but specific details provided to assist those of ordinary skill in the art in a comprehensive understanding of the invention, and the invention is only defined within the scope of the appended claims.
In the entire description of the invention, the same drawing reference numerals are used for the same elements across various figures. In the drawings, sizes and relative sizes of layers and areas may be exaggerated for clarity in explanation.
It will be understood that when an element or layer is referred to as being “on,” “connected to” or “coupled to” another element or layer, the element or layer can be directly on, connected or coupled to another element or layer or intervening elements or layers. In contrast, when an element is referred to as being “directly on,” “directly connected to” or “directly coupled to” another element or layer, there are no intervening elements or layers present. As used herein, connected may refer to elements being physically and/or electrically connected to each other. As used herein, the term “and/or” includes any and all combinations of one or more of the associated listed items.
Although the terms “first, second, and so forth” are used to describe diverse elements, components and/or sections, such elements, components and/or sections are not limited by the terms. The terms are used only to discriminate an element, component, or section from other elements, components, or sections. Accordingly, in the following description, a first element, first component, or first section may be a second element, second component, or second section.
The terminology used herein is for the purpose of describing particular exemplary embodiments only and is not intended to be limiting of the invention. As used herein, the singular forms “a,” “an” and “the” are intended to include the plural forms as well, unless the context clearly indicates otherwise. It will be further understood that the terms “comprises” and/or “comprising,” when used in this specification, specify the presence of stated components, steps, operations, and/or elements, but do not preclude the presence or addition of one or more other components, steps, operations, elements, and/or groups thereof.
Embodiments of the invention are described herein with reference to cross-section illustrations that are schematic illustrations of idealized embodiments (and intermediate structures) of the invention. As such, variations from the shapes of the illustrations as a result, for example, of manufacturing techniques and/or tolerances, are to be expected. Thus, embodiments of the invention should not be construed as limited to the particular shapes of regions illustrated herein but are to include deviations in shapes that result, for example, from manufacturing.
Unless otherwise defined, all terms (including technical and scientific terms) used herein have the same meaning as commonly understood by one of ordinary skill in the art to which this invention belongs. It will be further understood that terms, such as those defined in commonly used dictionaries, should be interpreted as having a meaning that is consistent with their meaning in the context of the relevant art and will not be interpreted in an idealized or overly formal sense unless expressly so defined herein. All methods described herein can be performed in a suitable order unless otherwise indicated herein or otherwise clearly contradicted by context. The use of any and all examples, or exemplary language (e.g., “such as”), is intended merely to better illustrate the invention and does not pose a limitation on the scope of the invention unless otherwise claimed. No language in the specification should be construed as indicating any non-claimed element as essential to the practice of the invention as used herein.
In the laser induced thermal imaging (“LITI”) method, a laser beam emitted from a laser source is patterned using a. mask having patterns. To form the transmissive regions and the non-transmissive regions of the mask, reflective patterns can be used. Generally, the reflective patterns are manufactured by a multistage process, that is, a photolithography process, including forming a light-reflecting layer on a mask substrate, coating photosensitive resin (otherwise referred to as photoresist) on the light-reflecting layer, and performing pattern exposure, development. etching using the resist as a mask, and the removal of the resist.
Therefore, this method of manufacturing a mask requires a multistage process. However, as the number of photolithography processes increases, process efficiency is reduced. In particular, this method is not suitable for manufacturing a large-area mask. Therefore, there remains a need for an improved mask and manufacturing method thereof, which has increased process efficiency and is suitable for large-area mask manufacturing.
Hereinafter, exemplary embodiments of the invention will be described with reference to the attached drawings.
FIG. 1
FIGS. 2 through 7
FIG. 1
is a flowchart illustrating an exemplary embodiment of a method of manufacturing a mask according the invention. are cross-sectional views illustrating exemplary embodiments of operations in the manufacturing method of .
FIG. 1
110
130
150
170
Referring to , the method of manufacturing the mask includes providing a base substrate having light-absorbing layer patterns on a first surface thereof (operation S), forming (e.g., providing) a reflective layer on the light-absorbing layer patterns and the base substrate (operation S), and providing reflective patterns by partially removing the reflective layer (operation S). The method may further include providing a reflection preventing layer on a second surface of the base substrate (operation S).
FIGS. 1 through 4
110
Referring to , the providing of the base substrate having the light-absorbing layer patterns on the first surface thereof (operation S) may be performed as follows.
FIG. 2
120
110
First, referring to , a light-absorbing material layer is formed on a first surface of a base substrate .
110
110
110
The base substrate may include a light-transmitting material that allows light such as laser light to transmit therethrough. However, there is no restriction on the material of the base substrate . In one exemplary embodiment, for example, the base substrate may include a glass substrate, a quartz substrate, a sapphire substrate, a ceramic substrate, a semiconductor substrate, etc.
120
120
120
The light-absorbing material layer may include a material that can absorb light such as laser light, such as a conductive material or an insulating material. The conductive material may be, but is not limited to, an element of chrome (Cr), molybdenum (Mo), nickel (Ni), titanium (Ti), cobalt (Co), copper (Cu) or aluminum (Al), an alloy material containing the above element as a main component thereof, a compound (such as a nitrogen compound, an oxygen compound, a carbon compound, or a halogen compound), etc. In one exemplary embodiment, when the conductive material is used, the light-absorbing material layer may be formed by a deposition method, a sputtering method, a chemical vapor deposition (“CVD”) method, etc. When the insulating material is used, the light-absorbing material layer may be formed by a coating method.
120
The light-absorbing material layer may also include a semiconductor material such as silicon germanium, molybdenum oxide, tin oxide, bismuth oxide, vanadium oxide, nickel oxide, zinc oxide, gallium arsenide, gallium nitride, indium oxide, indium phosphide, indium nitride, cadmium sulfide, cadmium telluride or strontium titanate; an organic resin material such as polyimide, acrylic, polyamide, polyimide-amide, resist or benzocyclobutene; or an insulating material such as siloxane or polysilazane.
FIG. 2
120
120
In , the light-absorbing material layer has a single layer structure. However, a single layer structure is merely an example, and the light-absorbing material layer can also have a multilayer stacked structure.
FIG. 3
131
120
131
131
131
Referring to , photoresist patterns are formed on the light-absorbing material layer . The formation of the photoresist patterns may be achieved by a known photolithography process. While a single photoresist pattern is labeled as in the figures, it will be understood that a plurality of photoresist patterns may also be represented by the label.
120
131
121
120
131
121
120
131
110
120
121
121
FIG. 4
The light-absorbing material layer is patterned using the photoresist patterns as a mask, thereby forming light-absorbing layer patterns as shown in . That is, portions of the light-absorbing material layer which are masked (e.g., overlapped) by the photoresist patterns may become the light-absorbing layer patterns , and portions of the light-absorbing material layer which are not masked by the photoresist patterns may be removed. In addition, the first surface of the base substrate may be exposed at areas where portions of the light-absorbing material layer have been removed. While a single light absorbing pattern is labeled as in the figures, it will be understood that a plurality of light absorbing patterns may also be represented by the label.
120
There is no restriction on the method of patterning the light-absorbing material layer . In one exemplary embodiment, for example, dry etching or wet etching can be used.
131
131
131
131
Although not shown in the drawings, the photoresist patterns may, if necessary, be removed before a subsequent process. Alternatively, the subsequent process may be performed without removing the photoresist patterns . In particular, if the subsequent process is performed without removing the photoresist patterns , a photoresist stripping process may be omitted, thus reducing the number of processes performed. An exemplary embodiment where the subsequent process is performed without removing the photoresist patterns will be described below as an example, but the invention is not limited thereto.
FIGS. 1 and 5
130
Referring to , the providing of the reflective layer on the light-absorbing layer patterns and the base substrate (operation S) may be performed as follows.
FIG. 5
FIG. 5
150
110
121
150
110
121
131
131
150
110
121
121
Referring to , a reflective material layer is formed on the first surface of the base substrate having the light-absorbing layer patterns thereon. That is, the reflective material layer is formed on the first surface of the base substrate exposed between the light-absorbing layer patterns and on the photoresist patterns . Although not shown in , if the photoresist patterns have been removed as described above, it follows that the reflective material layer can be formed on the first surface of the base substrate exposed between the light-absorbing layer patterns and on the light-absorbing layer patterns .
150
150
150
2
2
2
2
5
2
2
3
2
3
2
3
2
2
2
3
The reflective material layer may include a dielectric. The reflective material layer may include a dielectric material such as any one material or a compound of two or more materials selected from SiO, TiO, ZrO, TaO, HfO, AlO, ZnO, YO, BeO, MgO, PbO, WO, VOX, SiNX, eNX, MN, ZnS, CdS, SiC, SiCN, MgF, CaF, NaF, BaF, PbF, LiF, LaF, and GaP. However, the material of the reflective material layer is not limited to the above example materials.
150
150
150
150
150
150
150
150
150
40
150
150
150
150
150
a
c
b
d
a
d
FIG. 5
FIG. 5
FIG. 5
The reflective material layer may have a multilayer structure. In an exemplary embodiment, the reflective material layer may include first reflective films and and second reflective films and stacked alternately. In , the reflective material layer has a four-layer structure. However, the four-layer structure is merely an example used for ease of description, and there is no restriction on the number of layers included in the reflective material layer . In one exemplary embodiment, for example, the reflective material layer can have a structure consisting of or more layers. In , the first reflective film is located at the bottom of the reflective material layer as a lowest layer, and the second reflective film is located at the top of the reflective material layer as a uppermost layer. However, films of the lowest and uppermost layer shown in is merely an example, and reflective films located at the bottom and top of the reflective material layer can be changed if necessary.
150
150
150
150
150
150
150
150
150
150
150
150
150
150
150
a
c
b
d
a
c
b
d
a
c
b
d
2
2
3
2
Any one of the first reflective films and and the second reflective films and may include a high dielectric constant (e.g., high-k) material, and the other one of the first reflective films and and the second reflective films and may include a low-k material. In one exemplary embodiment, each of the first reflective films and may include a high-k layer such as TiOor AlO, and each of the second reflective films and may include a low dielectric constant (low-k) layer such as SiO. That is, the reflective material layer may have a structure in which the low-k layer and the high-k layer are stacked alternately. However, such a stacked structure is merely an example, and the material of the high-k layer and the material of the low-k layer can be selected from the above materials for forming the reflective material layer . In addition to the above materials for the reflective material layer , any of a number of materials that have been developed and commercialized or are realizable depending on future technological developments can be used as the material that forms the high-k layer and the material that forms the low-k layer.
150
150
150
150
150
150
a
c
a
c
a
c
The first reflective films and may include the same material, or any one of the first reflective films and may include a different material than the other. In one exemplary embodiment, for example, if the first reflective films and include a high-k material, the films may include the same high-k material or different high-k materials. Although not shown in the drawings, if the first reflective films include three or more layers, the layers may include the same material, or any of the layers may include a different material than another layer.
150
150
150
150
150
150
b
d
b
d
b
d
Likewise, the second reflective films and may include the same material, or any one of the second reflective films and may include a different material than the other. In one exemplary embodiment, for example, if the second reflective films and include a high-k material, the films may include the same high-k material or different high-k materials. Although not shown in the drawings, if the second reflective films include three or more layers, the layers may include the same material, or any of the layers may be formed of a different material than another layer.
150
150
There is no restriction on the method of forming the reflective material layer . The reflective material layer can be formed using any of a number of methods that have been developed and commercialized or are realizable depending on future technological developments, such as a spin coating method, a spray coating method, a screen printing method, an inkjet method, a dispensing method, a sputtering method, a physical vapor deposition (“PVD”) method, a CVD method, a plasma-enhanced chemical vapor deposition (“PECVD”) method, a thermal evaporation method, a thermal ion beam assisted deposition (“IBAD”) method, and an atomic layer deposition (“ALD”) method.
FIGS. 1 and 6
150
Referring to , the forming of the reflective patterns by partially removing the reflective layer (operation S) may be performed as follows.
FIG. 6
121
110
121
110
110
121
Referring to , the light-absorbing layer patterns are irradiated with laser light L. Here, a direction in which the laser light L is irradiated may be a direction from a second surface of the base substrate toward the light-absorbing layer patterns . The second surface of the base substrate refers to a surface opposite the first surface of the base substrate on which the light-absorbing layer patterns are located.
121
4
2
4
3
4
2
3
4
3
4
There is no restriction on the laser light L as long as the laser light L has energy absorbed by the light-absorbing layer patterns . In one exemplary embodiment, for example, the laser light L may be selected from ultraviolet laser light, visible laser light, and infrared laser light. As a laser oscillator capable of oscillating the laser light L, an excimer laser oscillator such as KrF, ArF or XeCl; a gas laser oscillator such as He, He—Cd, Ar, He—Ne or HF; a solid-state laser oscillator using a monocrystal of YAG, YVO, forsterite (MgSiO), YAlOor GdVOor a polycrystal (ceramic) of YAG, YO, YVO, YAIOor GdVOdoped with at least one of Nd, Yb, Cr, Ti, Ho, Er, Tm and Ta; or a semiconductor laser oscillator such as GaN, GaAs, GaAlAs or InGaAsP may be used. Further, a fundamental wave or any of second to fifth harmonics may be used in the solid-state laser oscillator.
−12
−15
In addition, a continuous-wave laser beam or a pulsed laser beam may be used as the laser light L. As for the pulsed laser beam, an oscillation frequency of several tens of hertz (Hz) to several kilohertz (kHz) is typically used. However, a pulsed laser capable of emitting a laser beam with an oscillation frequency of 10 megahertiz (MHz) or more (which is far higher than the typical oscillation frequency) and a pulse width on the order of picoseconds (10seconds) or on the order of femtoseconds (10seconds) may also be used.
A cross-sectional shape of the laser light L may be circular, ellipsoidal, rectangular or linear (specifically, shaped like a relatively long and thin rectangle).
121
121
The energy of the laser light L may be high enough to cause emission of gas from the light-absorbing layer patterns or evaporation of the light-absorbing layer patterns .
FIG. 6
FIG. 5
FIG. 5
FIG. 5
110
121
121
121
153
150
121
150
110
151
121
121
121
153
150
121
As shown in , the laser light L transmits through the base substrate and is absorbed by the light-absorbing layer patterns . Then, portions of the light-absorbing layer patterns at which the laser light L arrives are laser-ablated, and the light-absorbing layer patterns and portions of the reflective material layer (see ) which are disposed on the light-absorbing layer patterns are removed together. Accordingly, portions of the reflective material layer (see ) which remain on the base substrate form reflective patterns . With laser ablation, the laser irradiated portions of the light-absorbing layer patterns are evaporated by the energy of the laser light L absorbed by the light-absorbing layer patterns , and the light-absorbing layer patterns and the portions of the reflective material layer (see ) which are disposed on the light-absorbing layer patterns are removed (or scattered).
150
110
110
121
121
153
150
121
FIG. 5
FIG. 5
There is no restriction of the method of laser irradiation. In one exemplary embodiment, if the reflective material layer (see ) does not include a light-absorbing material, the laser light L may be irradiated to the entire second surface of the base substrate . When the laser light L may be irradiated to the entire second surface of the base substrate , since only the light-absorbing layer patterns absorb the laser light L, only the light-absorbing layer patterns and the portions of the reflective material layer (see ) which are disposed on the light-absorbing layer patterns may be removed.
150
110
150
121
151
FIG. 5
FIG. 5
Alternatively, the laser light L may be irradiated in a. state where a. photomask is placed. In one exemplary, if the reflective material layer (see ) includes light-absorbing material, a photomask may be placed on a rear surface of the base substrate in order to reduce or effectively prevent damage to the reflective material layer (see ) due to laser irradiation, and then the laser light L may be irradiated. In using the photomask, the photomask may include opening regions corresponding to the light-absorbing layer patterns and blocking regions corresponding to the reflective patterns which are to be formed.
121
121
Alternatively, laser oscillators may be placed at positions corresponding to the light-absorbing layer patterns , and only the light-absorbing layer patterns may be irradiated with the laser light L. That is, there is no restriction on the method of laser irradiation.
2
110
110
110
110
Although not shown in the drawings, after or during the irradiation of the laser light L, a gas such as Nor air may be sprayed onto the second surface of the base substrate irradiated with the laser light L. In addition, the base substrate may be cleaned using a liquid such as water which is a non-reactive material. By spraying a gas onto the base substrate or by cleaning the base substrate with a liquid, it is possible to remove dust, foreign matter, etc. that affect laser ablation.
151
151
151
110
151
Since reflective patterns are formed using laser ablation as described above, a photoresist coating process, an exposure process, a development process using a developing solution, and a photoresist stripping process can be omitted from the process of forming the reflective patterns . The operation using the laser ablation can reduce or effectively prevent loss of materials such as a photoresist material and a developing solution, and simplify a mask manufacturing process. Furthermore, since the reflective patterns are formed by laser irradiation, even if the base substrate has a relatively large planar area, the reflective patterns can be formed more easily. That is, according to the invention, a large-area mask can be manufactured more easily.
150
151
FIG. 5
In a conventional photolithography process, an exposure process using a metal mask is performed to etch the reflective material layer (see ). To this end, the metal mask should additionally be manufactured according to the reflective patterns that are to be formed. According to the invention, however, the additional process of manufacturing the metal mask can be omitted, thereby improving process efficiency.
FIGS. 1 and 7
170
Referring to , the method of manufacturing the mask according to the illustrated exemplary embodiment may further include the providing of the reflection preventing layer on the second surface of the base substrate (operation S).
FIG. 7
170
110
That is, referring to , a reflection preventing layer may be formed on the entire second surface of the base substrate .
170
170
170
170
170
170
170
a
b
FIG. 7
FIG. 7
The reflection preventing layer allows incident light to transmit therethrough and reduces or effectively prevents incident light from being reflected. The reflection preventing layer may have a single layer structure or may have a multilayer structure including first and second layers and as shown in . When the reflection preventing layer has a multilayer structure, there is no restriction on the number of layers included in the reflection preventing layer . In , the reflection preventing layer has a two-layer structure. However, this structure is merely an example.
170
170
170
170
2
2
2
2
5
2
2
3
2
3
2
3
2
2
2
3
There is no restriction on the material of the reflection preventing layer . In one exemplary embodiment, the reflection preventing layer may include a material that maintains stable optical properties even if exposed to light (e.g., laser light) for a long time and/or a dielectric material such as any one material or a compound of two or more materials selected from SiO, TiO, ZrO, TaO, HfO, AlO, ZnO, YO, BeO, MgO, PbO, WO, VOX, SiNX, eNX, MN, ZnS, CdS, SiC, SiCN, MgF, CaF, NaF, BaF, PbF, LiF, LaF, and GaP. However, the material of the reflection preventing layer is not limited to the above example materials. The reflection preventing layer may be formed using any of a number of methods that have been developed and commercialized or are realizable depending on future technological developments, such as a spin coating method, a spray coating method, a screen printing method, an inkjet method, a dispensing method, a sputtering method, a PVD method, a CVD method, a PECVD method, a thermal evaporation method, a thermal IBAD method, and an ALD method.
170
In the reflection preventing layer having a multilayer structure, a low-k layer and a high-k layer may be stacked alternately so as to reduce or effectively prevent reflection of light. There is no restriction on the stacking order of the high-k layer and the low-k layer, the material of each layer and the cross-sectional thickness of each layer as long as reflection of light can be reduced or effectively prevented.
170
110
170
According to the illustrated exemplary embodiment, the reflection preventing layer is formed on the second surface of the base substrate after the laser ablation process. Therefore, even if the laser light L with high energy is used during the laser ablation process, the probability of damage to the reflection preventing layer due to the laser light L can be reduced or effectively prevented.
FIG. 8
FIGS. 9 through 12
FIG. 8
is a flowchart illustrating another exemplary embodiment of a method of manufacturing a mask according to the invention. are cross-sectional views illustrating exemplary embodiments of operations in the manufacturing method of .
FIG. 8
FIGS. 1 through 7
210
230
250
Referring to , the method of manufacturing the mask includes providing a base substrate having light-absorbing layer patterns on a first surface thereof and a reflection preventing layer formed on a second surface thereof (operation S), forming (e.g., providing) a reflective layer on the light-absorbing layer patterns and the base substrate (operation S), and providing reflective patterns by partially removing the reflective layer (operation S). Hereinafter, a description of elements and features identical to those described above with reference to will be made briefly or omitted.
FIGS. 8 through 10
210
Referring to , the providing of the base substrate having the light-absorbing layer patterns on the first surface thereof and the reflection preventing layer provided on the second surface thereof (operation S) may be performed as follows.
FIG. 9
FIG. 9
FIG. 7
170
110
170
170
170
170
170
a
b
Referring to , a reflection preventing layer is formed on a second surface of a base substrate . The reflection preventing layer reduces or effectively prevents reflection of external light. The reflection preventing layer may have a single layer structure or have a multilayer structure including first and second layers and as shown in . Other features of the reflection preventing layer are the same as those described above with reference to , and thus a description thereof will be omitted.
FIG. 10
121
110
121
110
131
131
Referring to , light-absorbing layer patterns are formed on a first surface of the base substrate . The light-absorbing layer patterns may be formed by forming a light-absorbing material layer on the first surface of the base substrate , forming photoresist patterns on the light-absorbing layer, and then patterning the light-absorbing material layer using the photoresist patterns as a mask.
FIGS. 9 and 10
FIGS. 9 and 10
170
121
121
110
170
110
In , the reflection preventing layer is formed before the light-absorbing layer patterns are formed. However, the order of these operations is merely an example. That is, unlike the illustration in , the light-absorbing layer patterns may be formed on the first surface of the base substrate , and then the reflection preventing layer may be formed on the second surface of the base substrate .
131
131
FIG. 3
Although not shown in the drawings, the photoresist patterns may, if necessary, be removed before a subsequent process, or the subsequent process may be performed without removing the photoresist patterns as described above with reference to .
FIGS. 8 and 11
230
210
Referring to , the providing of the reflective layer on the light-absorbing layer patterns and the base substrate (operation S) may be performed as follows after the providing of the base substrate having the light-absorbing layer patterns formed on the first surface thereof and the reflection preventing layer formed on the second surface thereof (operation S).
FIG. 11
FIG. 11
150
110
121
150
110
121
131
131
150
110
121
121
Referring to , a reflective material layer is formed on the first surface of the base substrate having the light-absorbing layer patterns thereon. That is, the reflective material layer is formed on the first surface of the base substrate exposed between the light-absorbing layer patterns and on the photoresist patterns . Although not shown in , if the photoresist patterns have been removed, the reflective material layer may be formed on the first surface of the base substrate exposed between the light-absorbing layer patterns and on the light-absorbing layer patterns .
150
150
150
150
FIG. 5
The reflective material layer may be formed to have a multilayer structure, but the structure of the reflective material layer is not limited to the multilayer structure. The method of forming the reflective material layer and the material that forms the reflective material layer are as described above with reference to .
FIGS. 8 and 12
250
230
Referring to , the forming of the reflective patterns by partially removing the reflective layer (operation S) may be performed as follows after the forming of the reflective layer on the light-absorbing layer patterns and the base substrate (operation S).
FIG. 12
121
110
121
110
110
121
Referring to , the light-absorbing layer patterns are irradiated with laser light L. Here, a direction in which the laser light L is irradiated may be a direction from the second surface of the base substrate toward the light-absorbing layer patterns . The second surface of the base substrate refers to a surface opposite the first surface of the base substrate on which the light-absorbing layer patterns are located.
170
110
110
121
121
121
153
150
121
150
110
151
FIG. 5
FIG. 5
The laser light L transmits through the reflection preventing layer to reach the base substrate , transmits through the base substrate and is absorbed by the light-absorbing layer patterns . Then, portions of the light-absorbing layer patterns at which the laser light L arrives are laser-ablated, and the light-absorbing layer patterns and portions of the reflective material layer (see ) which are disposed on the light-absorbing layer patterns are removed together. Accordingly, portions of the reflective material layer (see ) which remain on the base substrate form reflective patterns .
121
121
170
FIG. 7
There is no restriction on the laser light L as long as the laser light L has energy absorbed by the light-absorbing layer patterns . In particular, the intensity of the laser light L may be adjusted within a range that allows the laser light L to have energy absorbed by the light-absorbing layer patterns and does not substantially affect the reflection preventing layer . In addition, a continuous-wave laser beam or a pulsed laser beam may be used as the laser light L. A cross-sectional shape of the laser light L may be circular, ellipsoidal, rectangular or linear (specifically, shaped like a relatively long and thin rectangle). There is no restriction on the method of laser irradiation as described above with reference to .
170
110
170
According to the illustrated exemplary embodiment, since the reflection preventing layer is provided on the base substrate which is formed in advance, the process of forming the reflection preventing layer and other processes can be performed in parallel.
FIGS. 1 through 7
151
In addition, an additional process of manufacturing a metal mask can be omitted as described above in the embodiment of . Thus, process efficiency can be improved. Further, since the reflective patterns are formed by laser irradiation, a relatively large-area mask can be manufactured more easily.
FIG. 13
FIGS. 14 and 15
FIG. 13
1000
1000
is a cross-sectional view of an exemplary embodiment of a mask according to an exemplary embodiment of the invention. are enlarged cross-sectional views of exemplary embodiment of portions of the mask shown in .
FIGS. 13 through 15
1000
110
151
110
170
110
Referring to , the mask includes a base substrate , reflective patterns disposed on a first surface of the base substrate , and a reflection preventing layer disposed on a second surface of the base substrate .
1000
1
151
2
151
1
2
The mask includes first regions P in which the reflective patterns are disposed and second regions P which are disposed between the reflective patterns separated from each other. The first regions P reflect incident light (e.g., laser light), thus functioning as light-blocking regions in a patterning process. The second regions P transmit incident light (e.g., laser light), thus functioning as light-transmitting regions in the patterning process.
110
1000
151
170
110
110
FIG. 2
The base substrate forms the body of the mask and supports the reflective patterns and the reflection preventing layer . The base substrate may include a light-transmitting material. In one exemplary embodiment, for example, the base substrate may include, but is not limited to, a glass substrate, a quartz substrate, a sapphire substrate, a ceramic substrate, a semiconductor substrate, etc., as described above with reference to .
151
151
151
151
151
151
151
151
151
151
151
150
150
150
150
a
c
b
d
a
c
b
d
a
c
b
d
FIG. 5
FIG. 5
FIG. 5
The reflective patterns reflect incident light (e.g., laser light). The reflective patterns may have a multilayer structure. In one exemplary embodiment, each of the reflective patterns may have a structure in which first reflective film patterns and and second reflective film patterns and may be stacked alternately. The first reflective film patterns and and the second reflective film patterns and are substantially the same as the first reflective films and (see ) and the second reflective films and (see ) described above with reference to , and thus a detailed description thereof will be omitted.
170
170
170
170
170
a
b
FIG. 13
The reflection preventing layer transmits incident light (e.g., laser light) and reduces or effectively prevents reflection of the incident light. The reflection preventing layer may have a single layer structure or have a multilayer structure including layers and as shown in . The reflection preventing layer may include a material that maintains stable optical properties even if exposed to light (e.g., laser light) for a relatively long time.
FIG. 14
FIG. 13
FIGS. 13 and 14
FIGS. 1 through 12
FIG. 14
151
151
151
1000
151
151
1511
is an enlarged view of an exemplary embodiment of a portion ‘A’ of a reflective pattern shown in , more specifically, a. side surface of the reflective pattern . Referring to , the reflective patterns of the mask according to the illustrated exemplary embodiment are formed by forming a reflective material layer as a single piece and then partially removing the reflective material layer using a laser ablation method, as described above with reference to . That is, the reflective patterns are provided by physically removing a portion of the reflective material layer without an etching process. Therefore, as shown in , side surfaces of the reflective patterns may become uneven (as indicated by reference numeral ) in the process of partially removing the reflective layer.
FIG. 15
FIG. 13
FIGS. 13 and 15
151
151
151
1000
151
is an enlarged view of an exemplary embodiment of a portion ‘B’ of a reflective pattern shown in , more specifically, an edge of a top surface of the reflective pattern . Referring to , the reflective patterns of the mask according to the illustrated exemplar embodiment are formed by forming a reflective material layer as a single piece and then physically removing a portion of the reflective material layer without an etching process. Therefore, an edge of the top of the side surface of each of the reflective patterns may be extended further than a remainder of the side surface due to the process of partially removing the reflective material layer.
1513
151
2
110
1513
151
1
110
151
151
1513
As a result, a protrusion may be formed on the edge of the top of the side surface of each of the reflective patterns . That is, a height D from the first surface of the base substrate to the edge (e.g., the protrusion ) of the top surface of each of the reflective patterns may be greater than a height D from the first surface of the base substrate to the top surface of each of the reflective patterns excluding the protruded edge. In other words, the height of each of the reflective patterns may have a maximum value at the protruded edge (e.g., the protrusion ) of the top surface thereof.
FIG. 13
151
1
151
2
1
151
2
Further, referring to , since the reflective patterns are formed by forming a reflective material layer as a single piece and then physically removing part of the reflective material layer without an etching process, a width W of a topmost portion of each of the reflective patterns may be different from a width W of a bottommost portion thereof In one exemplary embodiment, for example, the width W of the topmost portion of each of the reflective patterns may be greater or smaller than the width W of the bottommost portion thereof.
151
1000
1000
That is, since the reflective patterns of the exemplary embodiment of the mask according to the invention are formed by partially removing a reflective material layer using a laser ablation method, the mask according to the invention is structurally different from other masks having reflective patterns formed using a conventional lithography method.
One or more exemplary embodiment of the invention provides at least one of the following advantages.
Since reflective patterns are formed by laser ablation, a method of manufacturing a mask in a simplified process and with improved process efficiency can be provided.
In addition, the invention provides a mask manufactured with improved efficiency.
However, the effects of the invention are not restricted to the ones set forth herein. The above and other effects of the invention will become more apparent to one of daily skill in the art to which the invention pertains by referencing the claims.
While the invention has been particularly shown and described with reference to exemplary embodiments thereof, it will be understood by those of ordinary skill in the art that various changes in form and detail may be made therein without departing from the spirit and scope of the invention as defined by the following claims. The exemplary embodiments should be considered in a descriptive sense only and not for purposes of limitation.
BACKGROUND
SUMMARY
DETAILED DESCRIPTION
BRIEF DESCRIPTION OF THE DRAWINGS
The above and other features of the invention will become more apparent by describing in detail exemplary embodiments thereof with reference to the attached drawings, in which:
FIG. 1
is a flowchart illustrating an exemplary embodiment of a method of manufacturing a mask according to the invention;
FIGS. 2 through 7
FIG. 1
are cross-sectional views illustrating exemplary embodiments of operations in the manufacturing method of ;
FIG. 8
is a flowchart illustrating another exemplary embodiment a method of manufacturing a mask according to the invention;
FIGS. 9 through 12
FIG. 8
are cross-sectional views illustrating exemplary embodiments of operations in the manufacturing method of ;
FIG. 13
is a cross-sectional view of an exemplary embodiment of a mask according to the invention; and
FIGS. 14 and 15
FIG. 13
are enlarged cross-sectional views of exemplary embodiments of portions of the mask shown in . | |
This project was performed in two parts. The first was focused on light intensity as it affects performance. A randomized complete block design (RCBD) was performed. Broilers, Cobb 500 (n = 1584) were housed in 3 commercial houses (121.9 x 12.2 m). In each house birds were randomized and placed in 72 pens of 121.9 x 121.9 cm (22 bird/pen, males and females). All the treatment groups were provided with 24h light (L) during the first week and then 18L:6Dark (D) and 20 lux from day 7 to 14. The 3 intensity treatments of 5 lux (lx), 10 lx and 20 lx (24 replications) with 18L:6D were started at day 14 and continued until 40 days of age.
The second experiment was designed to determine if birds showed a preference for light intensity while eating. A RCBD was performed with 3 different light intensities. Cobb 500 broilers (n= 180), were housed in 1 commercial house. They were placed in 6 pens. Each pen had 3 rooms with a specific light intensity and one feeder so the birds could choose under which intensity to eat after 14d of age. Feed disappearance for each feeder was collected and the lighting program was the same as in trial number one. Also a camera was set to record the feeding behavior of the birds (number of birds per treatment during one hour at a random time during the daylight period, before light turns off and one hour after light turns on).
The results suggest that from a welfare perspective meat-type broiler chickens prefer to eat and drink under 20 lx rather than 5 lx which is the common commercial practice. Results suggest that a greater attention to light intensity, particularly with respect to feeder placement, may not only benefit production performance, but also bird welfare due to their preference for increased light intensity when feeding.
Raccoursier Frost, Maurice, "Effect of Light Intensity on Production Parameters and Feeding Behavior of Broilers" (2016). Theses and Dissertations. 1854. | https://scholarworks.uark.edu/etd/1854/ |
Evidence from archaeological sites provides evidence of the growth of legumes in ancient Egypt. White beans are some of the agricultural products produced for export in Egypt. Egypt white beans belong to the Fabaceae family along with other legumes. The botanical name for white beans (white kidney beans) is phaseolus Vulgaris.
White beans grow in well-aerated and drained soils with a PH of 6 to 6.8. The plant requires a lot of water and warm weather for a good yield.
There are two types of Egypt white beans; type I and type II. The two varieties are grown during different seasons. Type I is grown from December to April, which is usually the winter season. Type two grows during summer, that is, mid-May to October.
180-200 Egypt white beans weigh 100grams. Their humidity rate is at a maximum of 16%. In Egypt, the white beans are edible unripe as fruit. When dried out, its seeds are cooked by boiling. White kidney beans are also edible as vegetables in Egypt and other parts of the world, and its straws make good fodder for animals.
White beans are healthy and highly nutritious as they contain edible fibre, proteins, magnesium potassium and iron. The seeds also contain amino acids, fatty acids, calcium, vitamin C, Vitamin A, fats, sodium and Carbohydrates.
The health benefits of taking kidney beans are as follows.
They help in the fight for cancer through its antioxidant properties.
Its soluble and insoluble fibre content is good for digestion and balancing of cholesterol.
The folic acid content fight heart disease.
Vitamin B1 helps in improved cognitive ability
Finally, they help in detoxification of foods by removing sulfites commonly found in cooked meals.
Among the many advantages of white beans is that they can be stored for more extended periods when well dried in a dark environment without losing on quality. Like any other legumes, white beans fix atmospheric nitrogen into the soil increasing its fertility for other plants such as cereals.
For the edible pod beans, the collection is best done while they are young. When they start to ripen, they are best picked daily as the wait can result in toughening/hardening. The beans should also be harvested free of moisture to prevent diseases. When plucked, the seeds should be refrigerated unwashed in air-tight containers and only removed as they are being used. Chilled green beans can only stay fresh for up to four days. For more extended storage periods, the pods can be blanched, prickled or canned and frozen.
For dry beans, they should be harvested when they are dry and hard before they split. Upon collecting the seeds are removed from the pods. The seeds are then packaged and stored in a cool and dry place.
Among the biggest importers of white beans include Morocco, Romania, Bulgaria, Saudi Arabia, the USA Algeria and Tunisia among others.
You can place your order of Egypt white beans on this platform, and we will deliver it at your convenience.
Get Instant Quote
Are you a producer of Egypt White Beans or other products?
Apply to sell your produce directly to buyers worldwide. | https://www.selinawamucii.com/produce/cereals/egypt-white-beans/ |
An automatic rope ranging device for a winch, comprises a support (14), a transmission device and a rope ranging device, the transmission device consists of a pinion (4), an input transmission shaft (5), a conical gear mechanism (6), a worm and worm gear mechanism (7), a sliding gear (8), an output transmission shaft (10) and an input gear (9). the output transmission shaft (10) is connected with a right-and-left thread screw sleeve (16) by a key (23), the right-and-left thread screw sleeve (16) and a rope-ranging block (1) connected with a wire rope (15) form a special rolling screw pair structure by a rolling ball pressed-in device (18). The special rolling screw pair structure has the structure that the rolling balls (17) are arranged into a left row and a right row by left and right circular arc grooves (20) at two outside sides of a pressed-in end of the rolling ball pressed-in device (18) in a raceway (21) of the right-and-left thread screw sleeve (16), and the rolling balls are in point contact with two inner sides, the width (B) of the pressed-in end of the rolling ball pressed-in device (18) is slightly smaller than the width of the raceway (21) (delta 1), the height (H) of the pressed-in end is smaller than the depth of the raceway (21) (delta 2), the rolling balls are in point contact with the bottom of the raceway (21), the rolling ball pressed-in device (18) is pressed-in a opening hole (22) of the rope ranging block (13) and then is connected with the rope ranging block (13) by a screw (19) to form a whole. The automatic rope ranging device for the winch has simple structure, easy manufacture, reliable work, less wearing of the key components and long service life. | |
🕑 Reading time: 1 minute
The application of Civil engineering in the construction of Smart infrastructure is the foundation for all the other key elements in a Smart city like Smart property, Smart economy, Smart living, Smart governance and Smart environment. The underlying principle behind most of these elements is that they are well linked and that they generate data, which can be used to ensure the optimal use of resources and improve performance.
A precise and a unanimously accepted definition of Smart city has remained ambiguous till date. There are multiple interpretations available and a lot many amendments are being made to it every day, but there still remains a certain degree of imperfection in its definition. It is basically because of the wide range of components that come under the umbrella term- Smart City.
When tried to pick the mutual parts from all the accessible interpretations, a Smart city is the one that supervises and fuses conditions of all its critical infrastructures, including roads, bridges, tunnels, rails, subways, airports, seaports, communication, water, power, even major buildings, with Information and communication technologies (ICT), to better optimize its resources, plan its preventive maintenance activities, and monitor security aspects while maximizing services to its citizens.
The role of civil engineering in some key components of smart city infrastructure is discussed in this article.
Contents:
Smart Transportation
A smart city is all about all-round connectivity and accessibility. The storage and efficient transfer of data by and between roads, vehicles, highways, bridges, traffic lights and even the relevant buildings is capable of assisting with the public as well as commercial transportation management, route information system, vehicle control and safety, and traffic congestions.
Though the Intelligent transportation system has been there for a while, the latest generation of solutions comes with features like traffic prediction, analytics and decision support, traveller information, advisory services, ticketing and fare collection that are capable of enhancing the current system to a whole new level.
A lot of sensors need to be embedded into new and existing roadways, buildings, bridges, posts and signs that continuously gather data from passing vehicles. All the different kind of vehicles on the road need to be able to interact with each other and the infrastructure without any interruptions.
City planners, therefore, are required to go an extra mile and work with Civil engineers to upgrade their urban infrastructures to incorporate sensors and IoT devices.
Read More: What is Intelligent Transportation System? Its Working and Advantages
Smart Buildings
In a Smart city, Civil engineers have to plan the construction of a building giving due consideration to the installation of smart building management systems.
With the advent of IoT and AI, the working mechanisms for lighting systems, fire protection systems, security system, CCTV, HVAC have undergone a complete overhaul.
Today, motion sensors can sense when an area is vacant or occupied and turn off or on the lights and lower the temperature; crews can be monitored using passive RFID as they check in; IP networks let the users adjust the HVAC settings of the building over their smartphone or tablet so the desired temperature is achieved by the time they arrive.
Another important technology is the adoption of smart concrete sensors. These sensors are embedded within concrete at the time of placing and are able to relay necessary information like concrete's health, temperature etc. that ultimately helps with increased sustainability and infrastructure lifespan- a huge benefit to Civil engineers.
Read More: Smart Construction Materials – Applications in Civil Engineering
Smart Water and Irrigation Systems
IoT has revolutionized the conventional water and irrigation systems as well. Now, the system uses digital technology to help save water, reduce costs and increase the reliability and transparency of water distribution.
Physical pipe networks are embedded with sensors that help analyse the available flow and pressure data to determine aberrations (such as leaks) in real time to manage water flow. Real-time information can be provided to customers to help conserve water leading to lower water bills.
In the field of irrigation, Automatic irrigation system are set to replace the traditional irrigation systems. Again, the conjunction of IoT and sensors is used for controlling the basic switching mechanism of motors by sensing moisture present in the soil.
Moreover, Automatic irrigation system makes it possible to identify the amount of water to be delivered that will assist in maintaining the level of soil moisture and monitoring the level of water-tank which stores the water that will aid in the irrigation system.
Smart Sewerage Management
A smart sewerage management system is required to manage the flow of waste through low volume and high volume periods and, occasionally, at times of high precipitation that brings in a heavy influx of water, making the way into the sewers.
Smart sewerage systems allow city sewer infrastructure to store overflows in huge tanks (interceptors) constructed in various parts of the system. The addition of smart sensors to detect and monitor flow levels allows smart sewer systems to manage gates and valves that direct wastewater to locations where there is sufficient storage space.
Sensors are also used to monitor sewer lines for any weaknesses or damage that may require attention providing time to enable convenient scheduling of maintenance trips and routines. | https://theconstructor.org/building/civil-engineering-aspects-smart-city/33870/ |
Due to the fact that chemical reactions happen extremely quickly, it is difficult to identify the principles behind such reactions and the intermediary substances created during the process. To overcome these limitations and to predict and analyze chemical reactions, scientists depend on ‘computational chemistry’, which factors in the shape, movement, binding force or stability of the target particles into mathematical equations and then performs computer simulations.
The Computational Catalysis and Emerging Materials Laboratory headed by professor Jeong Woo Han at the Department of Chemical Engineering, POSTECH, deploys computational chemistry to design novel materials, predict their properties and validate them through experimentation with an aim to develop completely new catalysts and energy materials. Its research primarily focuses on ‘green’ catalysts used to produce fuel cell electrodes, store hydrogen or reduce the emission of exhaust gas.
One of the most notable research outcomes is related to the redesigning of platinum catalysts to improve their efficiency. These Pt catalysts, used to purify such exhaust gases as carbon monoxides or hydrocarbons into carbon dioxide and water, usually exist in clump formations made when the Pt atoms combine. The Pt atoms located on the inner side of these clumps are unable to participate in any chemical reaction. Leveraging computational chemistry, the Lab predicted that specifically-structured supports made from titanium and carbon ensure that Pt atoms do not form such clumps when separated. This was followed by physical experimentation to verify the hypothesis and improve the efficiency of Pt catalysts as a result. These research findings were featured on the 2019 January issue of the international journal of ‘ACS Energy Letters’, and the article made it on the MOST READ list for December 2018 as it was published in advance online.
In 2020, the Lab succeeded in developing an ultra-stable and highly active ceria catalyst used for CO oxidation. This catalyst was designed by co-doping rare earth and transition metals on ceria. Researchers first deployed computational chemistry to verify that simultaneously doping lanthanum (a rare earth metal) and copper (a transition metal) resulted in improvements in both activity and stability, and moved on to design a ceria catalyst. They proved through experimentation that this new type of ceria catalyst was as efficient as its conventional counterparts even at such a low temperature of 150 and yet, still remained stable for nearly 700 hours despite significant temperature fluctuations. The article was selected as a supplementary cover for the Dec. 2020 issue of the international journal of ‘ACS Catalysis’.
Computational chemistry gives researchers a head start as data that is accumulated each year enables them to predict the interactions among particles of interest with higher accuracy. It also allows them time and cost savings in their experimentations. Another plus is that the high-performance computers installed at the Lab can be connected to PCs or smartphones enabling researchers to work from any given location. The Lab aims to build a data network of energy materials to develop and commercialize novel catalysts and energy materials and to leverage this network to harness artificial intelligence in the development of catalysts and materials in order to establish methodologies that aid in the creation of never-before-possible materials. | https://home.postech.ac.kr/eng/computational-catalysis-and-emerging-materials-lab/?k&c&event |
Archive | Cost Accounting
This article throws light upon the seven main steps for installations of a costing system. The steps are: 1. Objectives to be Achieved 2. Study the Product 3. Study the Organisation 4. Deciding the Structure of Cost Accounts 5. Selecting the Cost Rates 6. Introduction of the System 7. A Follow-up. Installation of a Costing […]
This article throws light upon the three main types of costs. The types are: 1. Fixed Costs 2. Variable Costs 3. Semi-Variable Costs. Type # 1. Fixed Costs: Fixed Costs also referred to as non-variable costs, stand-by costs, period costs or capacity costs are those costs which do not vary with changes in volume of […]
Here is a compilation of top eight problems on break-even analysis with their relevant solutions. Break-Even Analysis: Problem with Solution # 1. From the following particulars, calculate: (i) Break-even point in terms of sales value and in units. (ii) Number of units that must be sold to earn a profit of Rs. 90,000. Solution: […]
After reading this article you will learn about Income Determination under Absorption and Marginal Costing. Under absorption costing, fixed costs are treated as product costs while marginal costing excludes fixed costs from product costs. The example given here illustrates the method of income determination under absorption and marginal costing: Example: In the two income statements […]
After reading this article you will learn about Differential Cost:- 1. Meaning of Differential Cost 2. Determination of Differential Cost 3. Essential Features 4. Managerial Applications. Meaning of Differential Cost: Differential costs are the increase or decrease in total costs that result from producing additional or fewer units or from the adoption of an alternative […]
Contribution is the difference between sales and variable cost or marginal cost of sales. It may also be defined as the excess of selling price over variable cost per unit. Contribution is also known as Contribution Margin or Gross Margin. Contribution being the excess of sales over variable cost is the amount that is contributed […]
This article throws light upon the ten major managerial problems in application of marginal costing. The problems are: 1. Pricing Decisions 2. Profit Planning and Maintaining a Desired Level of Profit 3. Make or Buy Decisions 4. Problem of Key 5. Selection of a Suitable or Profitable Sales Mix 6. Effect of Changes in Sales […]
Welcome to AccountingNotes.net! Our mission is to provide an online platform to help students to discuss anything and everything about Accounting. This website includes notes, procedures, problems and solutions and other allied information submitted by visitors like YOU.
Before publishing your Notes on this site, please read the following pages:
| |
ZMM Nova Zagora was established in 1970 as a main manufacturer of axles, shafts and gear wheels for all machine-building companies in Bulgaria at that time, belonging to a united state structure. Following its deep specialization in producing rotation parts, in 1974 the company started production of gear boxes for universal lathes and in 1977 of automatic gearboxes for CNC machines. In 1980 started production of assembly units for universal lathes.
In 1984 in the company included lathe CA 161 to its product list. The model was awarded a gold medal at the International Technical Fair Plovdiv 1986.
Over the years ZMM Nova Zagora became supplier of more than 30 machine-builders in the world. ZMM Nova Zagora has produced parts and details for more than 300,000 machine tools, gears for over 200 000 engines, over 15,000 meters automatic production lines for Germany, Russia, etc., and for the past 10 years over 200 000 piston rods for heavy hydraulic cylinders.
Since its establishment until now ZMM Nova Zagora has produced over 11 million gear ratchets (gears, shaft gears and spline shafts) embedded in various types of machine tools operating around the world.
The company became a part of ZMM Bulgaria Holding, together with IHB Metal Castings, ZMM Sliven and IHB Electric.
ZMM Bulgaria Holding is а part of Industrial Holding Bulgaria, the largest industrial group in Bulgaria. | http://en.zmmnz.com/about-us/history |
Farmer NGOs criticize Latvia’s green agricultural policy plans / Article
The total European Union investment for direct payments to Latvian farmers and rural development over the period 2021-2027 is around 3.4 billion euros. Although direct payments have increased by up to 44%, Latvian farmers will only see their funding increase by 11% due to the low level of state co-financing and the reduction in funds allocated to rural development.
Requirements and results
The Deputy State Secretary of the Ministry of Agriculture, Rigonda Krieviņa, said that according to the parameters of the European Union, agricultural aid in member states should be oriented towards environmental, economic and social objectives. Therefore, the last two years have produced a national strategic plan for the common agricultural policy.
“For measures of economic objectives, we have earmarked 42% of our total money, for environmental measures, they are 48% and for social purposes, which include education, science and innovation, we have planned 8.8% of funding. We have developed 66 measures in our strategic plan, together with the farmers’ association, ”Krieviņa said.
The strategy also provides for increased support for agricultural practices that contribute to the mitigation of climate change, the improvement of the environment and the preservation of biodiversity. 438 million euros are devoted to voluntary activities called eco-programs.
“When claiming additional payments, the farmer needs to know whether he will keep the water, soil, air clean or reduce CO2 emissions with his activities. [..] It will also benefit from very strict monitoring and the European Commission will monitor financial use. If these goals are not met, funding adjustments will also be made, ”Krieviņa said.
Association: The environmental benefits are insignificant
Gustava Norkārklis, head of the Latvian organic farming association, said supporting ecosystems, as planned to achieve Europe’s “green” goals, is more focused on strengthening large conventional farms than on development. from organic farming.
“It will benefit the environment very little, but a lot of money will be spent. I criticized these eco-programs. For example, one of the aid measures is that a farmer just has to perform soil tests and report on the fertilizer plans once every 5 years, and he will receive an aid payment for those years. . But this is just normal agricultural practice. I’m not sure what the significant improvements and benefits are for the environment and biodiversity. And there is significant funding – over 100 million, ”Norkārklis pointed out.
The European Commission “From farm to fork” The strategy aims to increase the share of organic farmland in the EU to 25% by 2030, while reducing the use of pesticides by 50%. Norkārklis estimated that funding for farmers to practice organic would remain at almost the same level, unlike the increase made by the ministry.
“The financing of organic farming is now divided into several components. In the past, there was only one payment for organic farming – a clear envelope and a budget. From now on, its money is divided between eco-programs and rural development. And it can be said that there is now another payment available for organic farmers, but it is the same old green payment. By replacing the title of an aid measure, it is not entirely correct to say that more funding has been given to organic farmers. We will encourage the European Commission to assess the strategic plan of our Common Agricultural Policy and indicate where we see concerns about it, ”said Norkārklis.
Inequalities between sectors
On the other hand, Farmers Saeima foreign policy specialist Valters Zelès said the organization representing mostly conventional farmers also sees many gaps in the national strategy, such as unequal aid payments across sectors. and production methods.
“The ministry has created a payment calculator where owners can enter their farm data individually and see how much aid will increase or decrease compared to the previous period. [..] Around 20% of farms will remain at the current support level and 20% could benefit from a 10-20% increase in support payments. But 50-60% of homeowners will experience a reduction of around 20% in the aid payment [continuing to operate as it has been].
“And the ministry also emphasizes that we cannot do things the way we have done so far. We must do more. But the most important problem is that the ministry did not predict that, for example, in agricultural production, all farmers could do more. Even if we wanted all farms to become so ‘green’ tomorrow and put all these eco-program practices in place, there would only be enough money for 25% of the farmers, ”Zelès said.
In the opinion of the “Farmers Saeima”, farmers in the organic sector will have an increase in payments, but the national plan for the common agricultural policy lacks targets to achieve.
“For example, there is a huge increase in payments in the biosector – over 100 million. But we don’t see how we can increase productivity or combat the fact that, for example, organic milk is produced, but transmitted like conventional milk. In particular, it appears that the consumer pays the farmer, through taxes, to produce the organic dairy, which results in mixing in one tank with conventional milk. Therefore, we do not see any strategy or logic in the way these payments are distributed. We don’t see an impact assessment done in Latvia for this whole strategic plan, ”Zelès said.
Focus on planning and reporting
Andrejs Briedis, director of the Latvian Nature Fund, believes that ecosystems developed by the ministry should have basic requirements to be able to apply for help, rather than providing them as innovative environmentally friendly solutions.
“At this point, we don’t have national targets, because what we can find out by reading the strategic plan is how much money and on what hectares we are going to spend. You can’t say that its eco-programs are all bad, but most of the time farmers don’t have to change anything or very little. In fact, the focus is on better planning and better accountability across the country for what farmers do.
But if we come back to the objectives of the Parcours Vert: to reduce the use of minerals and pesticides, then ecosystems do not foresee anything like this. So what strategy are we talking about here? Unfortunately, we have to conclude that the national strategy is to take all the money as much as possible, ”said Briedis.
The Environmental Advisory Council also analyzed the national strategy of the common agricultural policy and also concluded that the aid does not focus on more environmentally friendly agriculture. Council President Juris Jātnieks stressed that Latvia’s strategic plan does not increase funding or compensation for forest owners whose territories are located in specially protected natural areas.
Following the parameters of the new common agricultural policy strategy, support should be granted from 2023. The Ministry of Agriculture is currently developing proposals for public consultations and will submit the national plan to the European Commission for evaluation by now. the end of the year. Until the start of the new period in 2023, agriculture will benefit from transitional support conditions. | https://sari-organik.com/farmer-ngos-criticize-latvias-green-agricultural-policy-plans-article/ |
This project consisted of water main replacement on St. Paul Avenue from Benson Road to 39th Street North in Sioux Falls, South Dakota. Work included water main replacement, sanitary manhole repair, replacement of surfacing over trenches, and boulevard restoration. The water main replacement successfully utilized both open trenching and directional drilling methods. Sanitary manhole repair included reconstructing manhole benches and inverts along with installing manhole external frame seals.
This project consisted of installing 8” water main beginning approximately 1,700 feet southwest of the Great Bear entrance road and continuing 3,700 feet to the northeast. This project included water main installation, asphalt removal and replacement, and site restoration and grading. Also included in this project was the grading and paving for a turn lane approximately 600 feet long on Rice Street for the proposed TJN site.
This project consisted of installing approximately 4,500 feet of 16” DIP water transmission main on Six Mile Road between 10th Street and 26th Street in Sioux Falls, SD. This work included water main installation, directional boring 24” steel casing across SD Highway 42 and under a large box culvert, pavement replacement, and restoration.
This project consisted of approximately 2,200 feet of 6” and 8” water main replacement along with 1,400 feet of 8” sewer replacement. Also included with this project were street drainage improvements as well as the installation of ADA sidewalk ramps at 11th Street and Garfield Avenue.
This project consisted of 4,200 L.F. of street improvements on 15th Street (Clark Avenue to Garfield Avenue), Garfield Avenue (15th Street to 12th Street), and State Avenue (Thresher Drive to 12th Street). Also included in this project were 4,400 feet of water main, 2,000 feet of sanitary sewer, and two directionally drilled river crossings approximately 350 feet each. Street paving on 15th Street and Garfield Avenue was with 9″ concrete, with asphalt concrete paving used on State Avenue and at intersecting street tie-ins. ADA accessible sidewalks were also installed on this project for safe access to the middle school and high school. | https://dgr.com/tag/water-main/ |
The following comment was prepared by Dr N K Boardman, AO, FAA, FTSE, FRS for the meeting of the Prime Minister's Science and Engineering Council, 13 September 1996.
Knowledge, especially technological knowledge, is the main source of economic growth and improvement in the quality of life. Nations which develop and manage effectively their knowledge assets perform better(The OECD jobs strategy - technology, productivity and job creation, OECD, Paris, 1996)
The increasing recognition worldwide of the vital contribution of science and technology to economic growth, quality of life and environmental sustainability is reflected in the above quotation from the OECD. The OECD study on jobs strategy concluded that nations which develop and manage effectively their knowledge assets perform better.
The research which produces the knowledge base for innovation includes basic research, technological or applied research and engineering research. An important component of research in the universities is the training of students in research project planning and execution, creative thinking and initiative, and the use of up-to-date instrumentation and technology. Students who acquire these high-level skills are a vital resource for achieving a more competitive industrial sector, and improved environmental management and living standards.
The Australian Academy of Science emphasises the importance of high quality basic research and training to the knowledge base in science, which is referred to as the science base.
The science base includes fundamental research for the advancement of knowledge, often referred to as blue-sky research, and strategic basic research. The latter utilises the same techniques and instrumentation as fundamental research and is often long term, but it has the objective of contributing to definable problems.
The main research role of universities is the performance of fundamental research, although universities also perform much strategic basic research as well as some shorter term applied research. The balance between fundamental strategic basic and applied research is influenced by funding sources.
Medical research institutes perform both fundamental and strategic basic research with the aim of understanding and developing solutions for particular health problems. CSIRO and other government research agencies make a major contribution to the science base in their performance of strategic basic research which underpins and is essential to their shorter-term applied research.
The objectives of fundamental and strategic basic research of high international standard are to provide:
Published Australian research accounts for about 2% of the world's total. Effective international links are essential to provide access to the leading edge of world research and technology. Conducting high quality, internationally recognised research in Australia provides the entry ticket into the world community of researchers. Invitations to actively participate in international conferences and symposia and to visit leading research institutions overseas depend on the reputations of our scientists. International collaboration in research is also recognised as important for the visibility of Australian research.
A high quality science base needs:
The recruitment and retention of talented researchers are essential for the establishment and maintenance of research groups of high quality as benchmarked against world standards and performance. Employment conditions should be comparable to those in the USA and Europe. An issue for Australia is the proportion of researchers in tenured positions, compared with those on five-year renewable contracts or fixed-term appointments. Not everyone retains their research effectiveness for their whole working life, but in the universities the teaching, research and administrative loads can be varied.
Nurturing young talent in universities as well as in public research agencies is very important for the future health of the science base.
A good infrastructure for research, which includes laboratory space, equipment, libraries and computer facilities is essential for the performance of internationally competitive research and to attract talented individuals. A report (1) by the national Board of Employment, Education and Training (NBEET) concluded that research infrastructure, in all its dimensions is coming under increasing pressure due to expanding research activity in the universities.
A key problem in the university system is the unrealistic expectation that adequate research funding and research infrastructure will be available across all sections of higher education teaching. Should research resources particularly in some of the more expensive scientific disciplines be concentrated in fewer established institutions? Is critical mass important for the performance of internationally-competitive research?
The Academy of Science strongly supports the system of peer assessment of researchers and research projects and believes that the most gifted and able people should have adequate resources to enable them to perform research that is of high quality and significant in the international context. It is important that research resources also are provided for newly-appointed junior academic staff who, at the start of their career, often find it difficult to compete successfully with senior researchers with established track records.
The Academy of Science supports the continuation of a plurality of funding sources for research in universities while emphasising processes of peer assessment of research proposals and evaluation of performance and outcomes. There is concern, however, that the present success rate of 23% for applications for research grants from the large grant scheme of the Australian Research Council (ARC) is putting enormous pressure on the selection process and denying opportunities for many talented researchers.
A popular belief in the universities attributes the decline in research performance to inadequate funding so that researchers do not have adequate resources to undertake high quality internationally-competitive research.
In actual dollar terms (at constant prices) the total level of funding in the higher education sector has increased across all fields over the period 1984-1992. Across the period 1981-1990, there has been little change in the total research expenditure per research scientist and engineer indicating no change in the intensity of funding. This analysis does not take account, however, of the significant escalation in the costs of performing top quality internationally-competitive research in many areas of science during the same period.
Important elements in the strength and reputation of Australia's science base are centres of research excellence, such as the Research Schools and Centres of the Institute of Advanced Studies of the Australian National University, the ARC Special Research Centres and medical research Institutes. Such centres bring together a diversity of skills and a sufficient number of researchers and state-of-the-art equipment to tackle difficult research topics and mount long-term research programs in an internationally competitive way.
The report of the 1995 Joint Review (2) of the Institute of Advanced Studies of the ANU, commissioned jointly by the ARC and ANU, and the joint review reports of the individual schools and centres provide ample evidence for the high standing or research in the Institute of Advanced Studies and the important contributions to knowledge, and the science base in Australia, made by researchers at the Institute. The following quote is from the review report of the Institute as a whole.
'The Institute has acted as a magnet for talent. Its social, cultural and scientific environment has been such as to attract scholars of the highest calibre from all over Australia and indeed from all over the world. As a result the IAS is now a world player in every field in which it has well-established scholarly and research activity.'
Governments fund a large proportion of the research conducted in universities in most nations. For example, in 10 out of 14 countries of the OECD public money accounts for more than 80% of all funds going to university departments for research (3). The private sector, particularly in small and middle-size economies, is reluctant to fund long-term, blue-sky research and even strategic basic research because of the high risk and the inability of a company to appropriate the benefits of the research for competitive advantage. The results of university research are published in the open literature, although there is now a tendency in most universities to examine manuscripts for potentially valuable intellectual property before submission for publication. Even in the event of a delay in publication to enable patent application, the results become freely available to the scientific community as a knowledge base for further work and discovery.
In Australia, a substantial proportion of CSIRO's strategic basic research is long-term and high-risk and much of it is broadly applicable to a range of private sector activities. Research in areas of community interest such as the environment and public health is of increasing importance and clearly the responsibility of government.
An internationally accepted method of measuring the performance of basic research is by the quantity and quality of publications in international peer-reviewed journals.
Analysis of published Australian research papers in journals of the Science Citation Index as a share of world publications confirm the strength of Australia's basic research. Australia contributes 60% more papers to published science than its population or GDP would suggest (4).
Citations per paper in the international scientific literature are a measure of the visibility of the research and, with reservations, an indication of quality. Australia performs well across most scientific fields but with particular excellence in fields related to our national resources and competitive export industries: earth sciences, agriculture, plant and animal sciences and the environment. Australia also performs well in certain areas of medical research (5).
A disturbing feature of the analyses of Bourke and Butler (1993) and confirmed by the BIE (1996) is the declining share of world citations in a large number of fields of Australian research since the mid-to-late 1980s.
In a recent study, yet to be published, funded largely by the Australian Research Council, the Australian Academy of Science examined possible causes for the decline in citation share. Some evidence was obtained to support the view that the decline in the visibility of Australian science in the international scene is related to a reduction by Australian scientists in the tapping of international networks.
Overseas experience particularly at the post-doctoral level, or for PhD training, is an important way in establishing and maintaining networks. The proportion of academics in Australian universities who obtained their first degree in Australia and their PhD overseas has decreased from 21.5% in 1970 to 11.7% in 1994 (a decrease of 45%). Although it is difficult to obtain reliable statistics on post doctoral training overseas by Australian graduates, the opportunities have declined because sources of adequate overseas funding are more difficult to obtain. There are very few funding schemes in Australia to support overseas postdoctoral experience.
The Academy of Science believes that the lack of post-doctoral fellowships for study overseas is an important policy issue which has a bearing on future successful international collaborations.
The effective management of the knowledge assets of the science base requires a stronger R&D effort in the private sector and the establishment of stronger links between universities, government research agencies and industry. It is important not to put at risk the performance of the science base in universities and government research agencies by a confusion of research roles.
The public sector should not be coerced into doing research which is much better performed in the private sector because of its closeness to, and understanding of, markets.
The Academy of Science strongly supports the Cooperative Research Centres Scheme which is proving very successful in drawing together researchers from universities, government research agencies and industry and strengthening links with the users of research. Another objective of the scheme is a concentration of research resources for leading edge research in areas of national importance.
The Academy of Science puts forward the following as its priorities for the science base. | https://www.science.org.au/supporting-science/science-policy/submissions-government/comment%E2%80%94australian-science-and-technology%E2%80%93 |
For almost three decades now in SA's democratic project, we've managed as a nation to make formidable strides in various sectors of the economy. However, little to nothing, our efforts are showing tiny ripples of effect on a global scale mainly due to leadership poverty among other ills found in our governing systems. For most of us, we characterize our country as one that is reeling from poverty, structural inequality and gross unemployment as the deep challenges, which to an extent, is unacceptable given the material and immaterial resources our country has.
However, as is our reality with the elected leaders, our governing bodies have emblazoned the national psyche with the excruciating picture of the 'triple challenges'. We can thus argue that our scourges as a nation are mainly corollaries of the same thing: unstable and continuously weakening the economy. For many among ourselves, we're most often than not engaged in round-table dialogues convened by the both public and private domains and these are slowly becoming not only tiring, more, proving to be wasteful because they seldom bear any fruits for mutual benefit.
Hence the need to create more apolitical platforms for meaningful engagement. Also, we can attest that, given our current education curriculum, our system isn't responding or worse still, tackling present-future challenges of our national interest, but perpetuating the poverty plan inherited from our forebears. Like most cadres have said, we've heard more than enough complaints as a nation already, there's a huge need for innovative ideas to grow the socio-economy towards creating sustainable jobs.
Rationally, we may argue that sustainable jobs are not created by money per se, but by people through their innovative entrepreneurship. Given our status quo on global ranking, we may argue that on the one hand, our nation is replete with excellent research that does not lead to innovation and sustainable livelihoods; however, on the inverse, it is replete with 'entrepreneurs' that bring no sustainable innovation to the table but exists merely for the sake of public domain profiling. As a result, we have more tenderpreneurs than innovators.
A general counter held up to our national current economic predicament is that if our economy could be more robust, the number and quality of jobs should nominally improve. This would address not only our unemployment challenge but also lead to prospects of improved job spaces. In South Africa, the traditional mission of learning institutions has for a long while been to educate for employment by others rather than educating to innovate and, through entrepreneurship, employ others. Unfortunately, this mission, has for decades, been working against the ideals and aspirations of the time. Generation X and Y are fortunate to have platforms for expression.
When faced with a test and opponent of their presumed worth, most of our learning institutions have historically tended to claim that their role ends with education and the imparting of knowledge. Anything beyond that would thus be seen as a bonus. Practically, universities of technology in South Africa, in terms of their primary mandate, are supposed to focus on what is readily applicable in the workplace, by producing graduates who are work ready. In this context, these institutions have been called upon to strengthen their provisioning of work integrated learning, through which students are given an opportunity to experience first-hand the workplace as part of embedding their classroom learning.
Higher learning institutions should redirect their educational offerings towards deliberately cultivating a culture and ethos of innovative entrepreneurship across faculties.
Lately, we've learnt about a commendable initiative by government partnering with the private sector on making an increase of work integrated learning placements which represents a positive welcoming step in improving the educational efficacy in the country. Practically, this should generally strengthen the employment prospects of South Africans. Inversely, we remain hopeful that the move will bring a paradigm shift with positive impact towards the general economy. As a result, the work readiness for average job-seekers versus the economic potential of production will work.
As we may be aware that the world over, economies have for a long time struggled to produce adequate jobs to absorb skilled laborers, we can be hopeful that tectonic changes happening the world over will be embraced with the hope that most people, regardless of their formal training, will be absorbed by such large corporations. Ultimately, we need to acknowledge that the mechanism initiated by the public sector will combat our national challenges if we embrace the work integrated learning outcomes as necessary channels towards the realization of our NDP 2030 Vision.
In preparation for a brighter future
For a short while, SA's round-tables have become abuzz with the fancy word: innovation. At foundation phase, this modern concept has become part of our daily educational system at institutions embracing modern economy, but for the sake of inclusive economies, the term refers to creating new things: products and services. Ultimately, just like our international alliances at BRICS, our education system needs to re-orientate our people to be agents and interpreters of change, thus literally making them adept at noticing the myriad job opportunities that emerge from such changes.
Mindful of the fact that our higher learning institutions have made insignificant changes in teaching and learning, research and innovation, thus leading to almost irrelevant graduates to the workplace, it's equally imperative for the proper implementation of the work integrated learning to ensure the utilization of fresh graduates to avoid another generation of brain-drain. In short, we believe our higher learning institutions should redirect their educational offerings towards deliberately cultivating a culture and ethos of innovative entrepreneurship across faculties.
Dare we say that this call is made against the backdrop of sufficient international studies and models that show a positive correlation between innovative entrepreneurship versus economic prosperity. Mindful of our national colour defects and class struggles, we'll need to substantially re-orientate our people on the faculties at higher learning institutions to tamper with consideration of the resources and also the emotional fatigue it will carry. Ultimately and most importantly, we need to start somewhere to achieve our NDP 2030 Vision.
To fast-track progress, we envision going beyond some of the already existing intervention systems, but would most likely incorporate aspects of these as building blocks towards strengthening what's already in place. Our view is that a sustainable version of the national future requires not just the so-called entrepreneurship skills and opportunism that most people are chanting, but an innovation base from which to launch entrepreneurship.
SA's wide challenges in education and the economy cannot be solved through a single bullet.
To date, as we interact with graduates from a different background, it's evident that our higher learning institutions are grappling with how to effectively capitalize on some of their resourceful people and networks. This, however, can be alleviated by the consideration, incorporation and implementation of innovative approaches such as: leveraging partnerships with academic entrepreneurs -- entrepreneurs who have an affinity to academe -- and the work of entrepreneurial academics -- academics in higher learning institutions who are not introvert to publicly engage with industry.
Convinced of the urgency to move from where we are now as a civil society, this proposition, however, shouldn't be confused for a presumption that SA's wide challenges in education and the economy can be solely solved through a single bullet. All stakeholders must be aboard to achieve this common and shared future. | |
Science has the phenomenal potential to solve the world’s greatest challenges, transform businesses and make world a better place. From manufacturing personal, protective equipment, developing vaccines to the overall response to COVID-19, science is heightening the world’s expectations for what is possible, according to the 3M State of Science Index 2021.
The State of Science Index (SOSI) is a third-party, independent research study commissioned by 3M and conducted annually for the past four years to track attitudes towards science. The 2021 findings emphasize that there has been a significant improvement in the public image of science.
90% Indians trust science:
In India, 90% of respondents said that they trust science – a significant increase of 3 percentage points since the 2020 Pre-Pandemic Survey.
About 85% agreed that there are negative consequences to society if science is not valued while a majority of 87% Indian respondents will stand up to skeptics by defending science if someone questions it as against 75% globally.
Renewed interest in STEM careers; increased emphasis on gender diversity and inclusion in the field
While Covid-19 has certainly spotlighted the importance of scientists searching for a vaccine, its influence goes wider. The pandemic has ignited a renewed interest in STEM careers and education. Scientists and medical professionals are inspiring people to pursue STEM-based careers in the future, especially among younger generations.
In India, 91% of respondents agree that the world needs more people pursuing STEM-related careers. Due to the pandemic, about four-in-five are more inspired to pursue a STEM career that accounts for 83% as against 60% globally.
Inclusion in STEM
The survey also highlights the need for greater gender diversity and inclusion in STEM. About 85% of Indian respondents agree that it is important to increase diversity and inclusion in STEM fields and 83% acknowledge that underrepresented minority groups often do not receive equal access to STEM education.
“Events of the past year have put a spotlight on the education gap within underserved communities,” said Dr. Jayshree Seth, Corporate Scientist and Chief Science Advocate, 3M.
“Gender inequalities, and unequal access to a quality STEM education for under-resourced students, continue to affect economic outcomes across the globe. We must all do our part to create greater opportunities, by strengthening STEM investments, eliminating underrepresentation in STEM, and bridging the STEM talent gap so that we can all realize the promise of a more diverse, equitable, and inclusive society.”
“Science is becoming more of a uniting factor as the world moves toward a common mission to build a safer, greener, stronger, and more equitable future,” said 3M Chairman of the Board and Chief Executive Officer Mike Roman.
“The world’s confidence in science is confirmed every day as we see more and more examples of its impact, from the COVID-19 recovery to advancing sustainability, making a meaningful difference,” Mike Roman.
Need for immediate attention towards the health of the planet:
The State of Science Index exposes a growing concern and a sense of urgency surrounding the health of the planet. India has become more environmentally conscious – even more than many of the other countries surveyed.
A significant majority (87%) agree that solutions to mitigate climate change need to happen immediately, and 89% confirm their belief that the world should follow science to help create a more sustainable future.
Commenting on SOSI, 3M India Managing Director Ramesh Ramadurai said at 3M, the success of business is inextricably linked with the health of our planet.
“The pandemic has also brought greater attention to sustainability issues, with a focus on making more sustainable lifestyle choices. There is a lot one can do in our daily lives to bend the curve of environmental degradation. It is time we all realize the impact our action can make to change our environment,” said Ramesh Ramadurai.
Partner, preempt, prepare and prosper
The survey further emphasized the importance of cross-border and cross-sector collaboration as essential to scientific advancement in India. Approximately 84% of respondents feel that countries should collaborate to create solutions based on science to address major challenges.
Given the events over the past six months, the top three actions people in India most want corporations to prioritize are, prepare for future pandemics (52%), provide existing employees with new skills and training for their future careers (47%), and create new jobs/employment opportunities for underrepresented minority groups within their corporation (46%).
3M initiatives
3M recently announced actions to build even greater equity in its communities, business practices and workplaces, setting a new global, education-focused goal. The company will advance economic equity by creating five million unique STEM and Skilled Trades learning experiences for underrepresented individuals by the end of 2025.
3M is also releasing a docuseries for the public this June. “Not the Science Type,” features the stories of four female scientists with different careers as they challenge stereotypes and confront and overcome gender, racial, and age discrimination.
3M India Limited, an Indian Subsidiary of US based 3M Company, was established in 1988 and has its headquarters in Bengaluru with branch offices at Mumbai, Gurgaon, Pune, Kolkata and Chennai. 3M leverages its global innovation expertise to develop homegrown solutions that addresses unique needs of diverse customers in India.
3M has invested in Innovation centers at Bengaluru and Gurgaon to boost local product development. Its manufacturing footprint is spread across Bengaluru, Pune and Ahmedabad. From products that improve manufacturing efficiency and impact improved healthcare delivery to safety solutions that help increase road visibility, everyday kitchen aids and car care products, today, 3M science is improving the lives of millions of Indians. | https://thenfapost.com/2021/06/22/indians-bet-on-science-for-recovery-post-pandemic-reveals-3ms-state-of-science-index-sosi/ |
All returns must have a written RMA (Return Merchandise Authorization) number and must be reported to our team within 5 days of product receipt.
All returned merchandise must be unused, unwashed and in original packing.
The total return shipping cost will be paid by the customer.
A 20% Restocking fee will be charged.
Refunds will not be issued until the product is received and inspected by our team. | https://zuzulinens.com/pages/zuzu-supplies-shipping-returns-policy-usa |
FIELD OF THE INVENTION
The present invention relates to a self-appliable and removable lid lock for a conventional residential wheeled trash container. More particularly, the present invention relates to an automatic gravity inforced device to impove the means of locking and unlocking the hinged lid of a wheeled trash container. An advantage of the present disclosure of the self-appliable and removable lock and unlock device is not only that it can be purchased by the general public and be self attached effortlessly to a wheeled trash container, but can help avoid damage to a commercially owned receptacle due to complicated techniques other types of gravity inforced mechanisms require to drill holes in the plastic for mounting. A user can manually unlock the present invention by pulling an arm with a counterweight outwardly therefrom the device when attached to said trash container to disengage the mechanism and/or as the device is attached to said trash container standing in an upright position, is able to release the hinged lid from a locked position using the force of gravity as it tilts forward to be emptied by a waste service collection truck.
BACKGROUND OF THE INVENTION
As is known, conventional residential wheeled trash containers typically are invaded by scores of unwanted pests which selectively rummage through household bags of refuse at will because of insufficiient measures to keep the hinged lids closed, and with regard to the forgoing, giving free access also to unauthorized individuals to fill them at the costly expense to the customer. As another comparison, it is common for a customer to roll a wheeled trash container for placement in front of their residence for collection only to be struck by the hinged lid when strong winds throw it back suddenly without warning. This may be dangerous and cause injury, depending on the circumstance to the user. Most automated public utility garbage truck companies usually do not provide a means to their customers of locking the smaller wheeled trash containers for obvious reasons due to the finanical impact to supply every customer and the time involved for the drivers to manually reset each pickup if certain mechanisms used are not depenable to work properly on collection day.
DESCRIPTION OF THE RELATED ART
The related art of interests describes various devices for locking the hinged lids of residental trash containers that are permanently attached, intended to be unlocked when emptied by tilting into a waste collection service truck and then returned to a locked, upright position, although none discloses the present invention. There is a need for a dependable gravity inforced locking/unlocking device which permits instantaneous attachment for use to most conventional wheeled trash containers.
The related art will be discussed in order of perceived relevance to the present invention. These patents neither suggest nor do they teach the advantage of a hook and strap-on, lock and unlock mechanism offered for purchase to the general public to further improve an easily attachable and detachable means without damage to a waste receptacle by other examples shown.
As is known, automated public utility garbage truck companies do not supply a means of a lid lock device to the general public to solve the problem for unauthorized parties to gain access to the smaller waste collection containers waiting to be emptied. Conventional means to typically keep a smaller wheeled trash container hinged lid closed is frustrating for the general public at times and occanionally correct the problem by drilling holes into the plastic container body to attach a make-shift contrivance to hold the lid shut.
Moreover, many locking devices for locking the lids of trash containers have been proposed to solve these problems set forth in the specifications of claims herein, however, most locking devices affixed to the container rentals for public useage are applied by means of drilling holes into plastic receptacles solely for the purpose of attachment and are gravely unexceptable as indicated by the garbage truck companies due to unnecessary damage to their property. It should be noted, that a monotary value could be placed on the customer serving an excellent reason to deviate from this standard. It will be understood that an attempt to render a broader scope of a more practical application has been considered in the present invention as it is directly attached by means of hook and strap capability. In basic concept, the related art of interest, intended for ownership specifications to the public, can manually attach quickly without damage or prolonged wear and tear to the receptacle, perform its primary function and not interfere with the waste collection operation process.
Many locking devices for wheeled trash containers range in numbers to solve the problem of entry. For example, French Patent No. 2,721,912 to the assignee describes a locking device comprising a pivoting part pivotably mounted on a large, rectangular container referred to as “dumpster” body inside a protective casing. When said dumpster is upright, this pivoting part assumes a locked position in which it keeps the lid in a closed position. When the container is tilted for emptying, the pivoting part, by force of gravity, moves from its locked position to an unlocked position. Due to a dumpster's typically bulky size, the mechanism is generally directed to company owned containers that can be attached permanently.
The dumpster shown in U.S. Pat. No. 5,149,153 shows a huge container with hinged lids that solved the unauthorized use problem, by again, permanently attaching a self-disengaging lock between the hinged lid and the body of the dumpster, though, a key is required to release a latch. However, when locked, the latch is automatically released as the dumpster is tilted forward.
Although the smaller, residential trash containers have less problems with unauthorized entry which makes the problem seem unwarranted, smaller trash containers are prone to opening in high winds and animals have easier access. Refer to U.S. Pat. No. 5,738,395 that discloses a self-releasing latch arrangement applied to a coventional trash container having a hinged lid and being liftable and dumpable by the usual garbage truck. The latch arrangement comprises one or more keeper members, and once again, is permanently attached to the container outwardly thereof and overhanging thereform manually swinging the weight away from the keeper thereby emptying the container.
For more examples of prior art of locking devices for the hinged lid of a conventional trash container, especially gravity inforced mechanism attachments include U.S. Pat. No. 5,772,264, that discloses a gravity operated latch hook built in such a way as to incorporate a counterweight design by means of a hinge-plate and pin assembly, permitting the lid to open disengaging a striker when the container is set upright. Once again, the example of prior art is attached to the underside of the container's lid that requirers machine drilled holes into the receptacle with nut and bolt assembly.
U.S. Pat. No. 5,224,744 discloses a locking piece mounted to pivot between a position assuring the locking and unlocking position releasing the cover. The operating member is a movable weight mounted for guided sliding movement arranged so that a translatory movement in one direction of the operating member produces a rotation of the locking piece in the unlocking direction. Consequently, the device again is fabricated onto the receptacle.
While the embodiments of U.S. Pat. No. 5,772,061, claim a catch that includes a sliding/rolling locking member which is also the gravity element as indicated, the method of attachment connectable to the container requires complete assembly.
The disadvantages of the prior art designs disclosed in the patents listed herein are not only the difficulties of the lack thereof to the general public to attach such assembilies onto a commerically owned small residential trash containers quickly within seconds without the need of using any tools to drill holes for nuts and bolts or cutting edge welded applications, even though these devices fit most models of small wheeled trash containers with hinged lids on the market today, nor even though they may require some expertise for complexity and adjustment to assemble them when applied, but It must be understood that while the preferred methods seen possible, the most important factor to the equation is that if trash containers are not privately owned permission to do so must be granted by the commerical disposal service company.
However, none provide a hook and strap capability to facilitate a user to manually apply said device sufficiently and instantanously to a receptacle without damage to it or provide a removable locking arm system disposed adjacent to the container lid to prevent its opening. The present invention, in comparison, comprises a leverage capability locking means by when a lower pivoting arm with a counterweight rests in a vertical position is matched vertically to a connecting link plate. This dramaticly blocks movement of the second arm that is also connected by the opposite end of the connecting link plate in a horizontally locked orientation, preventing lift, serving as a temporary lock over the closed lid of the receptacle. As the container is tilted forward by a predetermined angle, the weighted first arm moves therefrom the housing structure, exerted by gravity, pushes against the link plate and dislodges the second arm to enable movement from its locked position, clearing the path of travel for the container lid for the waste contents to be removed. Having a properly located stop post on the housing structure to control limited movement for the lock arm and the counterweight arm is essential.
It can be seen that the present invention is an improvement over these prior-art locking devices for wheeled trash containers with hingable lids for attachment assemblies providing a simple, easy to apply without damage, cost efficient, fastenable to said trash containers with hingable lids and the like which meets the criteria desired for such devices.
SUMMARY
The object of the present invention is to further impove the means of locking and unlocking the hinged lid of a wheeled trash container, of the type comprising a fastenable housing structure on said container with strap, hooks and two rigid pivoting arms connected by a link plate; one horizonal and the other vertical. When said device has been properly hung on the outside edge of said container and the hinged lid closed, the upper horizonal pivoting arm holds the lid down in a locked position. As the receptacle is tilted forward, lifted and turned upside down for dumping this reaction triggers the lower vertical arm with a counterweight to swing away from the housing structure releasing the upper horizonal pivoting arm to move away from the lid and give clearence for dumping.
An illustrative embodiment of the wheeled trash container for the means of locking and unlocking the hinged lid using the force of gravity with the desire to prevent unauthorized access when the trash container is on the ground in an upright position without interference during the dumping process, discloses a device comprising a housing structure that directly engages the upper rim of a trash container below its hinged lid, by an upper arm whose purpose is to rest on top of the receptacle lid when in a closed position, for a locking means to prevent the lid from opening; and a lower arm with a counterweight both spaced-apart, each with a single affixed pivot point connected respectively by a movable link plate that directly serves as a locking or unlocking means when the container hinged lid is open or closed. In addition, when the trash container is in an upright position and waiting to be emptied, the key unlocking means of the device happens when the lower pivoting arm that rests in a vertical position against the receptacle pushes the link plate therefrom the trash container while it is tilted forward using the force exerted by gravity to raise the horizonal position of the upper arm to a vertical one above the receptacle. Furthermore, both pivoting arms comprise two stop tabs that have the added advantage of reducing movement when device is geometrical ingaged.
90
In further embodiment, the disclosure describes a pair of hooks on the housing structure that can attach the device to a trash container by fastening it under the hinged lid onto the edge of the outside wall. As a variant, said device may include an abutment acting as a retaining wall to resist lateral movement as well as a means capable of both keeping the mechanism in a upright position or keeping adequate distance therefrom while in direct contact with the container body. Thus, according to the invention, when it is hooked on the lip edge and tightly mounted with particular ease by a strap to the trash container's front lower handle, it can readily be installed on any existing container. When the trash container is tilted to a degree angle for dumping, the first arm with the counter-weight pivots outward by gravity pushing against the link and forcing the second arm from a locked position to an unlocked position giving clearence to the trash container's hinged lid for waste collection. When this object is accomplished because the trash container has been emptied and placed back on the ground after dumping, and the hinged lid is closed, the device returns to facilitate a locked position.
With regard to the foregoing, one embodiment of the disclosure provides a inexpensive locking mechanism for the general public to purchase and use without permanent alteration or damage to a privately owned commerical disposal service's property.
In a further embodiment, the disclosure describes a simple operation of a gravity operated device for locking a hinged lid to a trash container.
An advantage of the present disclosure may be that the locking device is able to automatically unlock during the process of dumping to eliminate any need for a user to manually unlock and lock during waste collection.
Additional objects and advantages of the disclosed device will be set forth in part in the description which follows, and/or can be learned by practice of the disclosed embodiments. The objects and advantages of the disclosure will be realized and attained by means of the elements and combinations particularly pointed out in the appended claims.
It is to be understood that both the foregoing general description and the following detailed description are exemplary and explanatory only and are not restrictive of the disclosure, as claimed.
BRIEF DISCRIPTION OF THE DRAWINGS
The invention will now be described, by way of example, with reference to the accompanying drawings, in which:
FIG. 1
is a perspective view of an illustrative embodiment of the automatic gravity inforced wheeled trash container hinged lid locking device;
FIG. 2
is a left side view of an illustrative embodiment of the automatic gravity inforced wheeled trash container hinged lid locking device showing the device in a locked position;
FIG. 3
is a right side view of an illustrative embodiment of the automatic gravity inforced wheeled trash container hinged lid locking device in a locked position;
FIG. 4
is a right side view of an illustrative embodiment of the automatic gravity inforced wheeled trash container hinged lid locking device showing the device in an unlocked position;
FIG. 5
is an exploded, front view of an illustrative embodiment of the automatic gravity inforced wheeled trash container hinged lid locking device;
FIG. 6
is a perpective view of an illustrative embodiment of the automatic gravity inforced wheeled trash container hinged lid locking device, showing the residential wheeled trash container equipped with a locking device standing in a closed upright position awaiting pickup by a garbage truck; and
FIG. 7
FIG. 1
illusrtates the operation of the device shown in , when the container is emptied shown with respect to a trash container open and an intermediate portion of the container and lid parts removed for convenience of illustration.
DETAILED DESCIPTION OF THE PREFERRED EMBODIMENT
FIGS. 1-7
1
3
7
8
4
5
6
10
30
9
1
1
2
41
5
6
4
37
Referring initially to of the drawings, an illustrative embodiment of the automatic gravity inforced trash container lid locking device shows a housing structure with hooks and and a pair of pivotable affixed arms and connected together by a movable link plate with properly designated stopper posts and and a strap adjustment chain . In accordance to the invention , an automatically lock and releasable holding arrangement , is interposed between lower container body and container hinged lid being held down in a locked position by a horizonal second arm and a cooperative link plate controlled by a vertical first arm with a weighted member .
1
It is to be understood that relative terms used herein are intended to describe relative relationships of components with respect to each other and should not be construed in a limiting sense. Therefore, the components designated by such terms may be oriented in other spatial relationships with respect to each other during use of the automatic gravity inforced trash container lid locking device .
1
3
4
5
6
7
8
10
30
37
9
7
8
43
2
2
FIG. 1
The automatic gravity inforced trash container lid locking device shown in comprises a housing embodiment having two pivoted arms and connected by a link plate , two hooks and , two stopper posts and , a counter weight , and an adjustment chain to automatically be attached with said hooks and and lashed to the front handle of a conventional automated trash container supplied by a dumping service to periododically lift and tilt the entire container to a upside down position for dumping waste in the usual manner.
3
10
4
37
6
5
4
5
3
7
8
42
2
3
1
2
40
1
2
FIG. 2
An illustrative embodiment comprises a housing structure shown in having a properly located stopper post for controllable movement of a first arm to facilitate a secured counterweight on one end, a pivotal location and movable link plate attached to a respective opposite end to a second arm with a pivotal location to further connect the first arm and the second arm to each other. The housing structure further comprises two selective hooks and for the purpose of hanging on the outside edge of an open conventional trash container , while the housing structure has an abutment that keeps the device in a somewhat vertical transition to the receptacle and presents a means of fastening a type strap or restraint from the device securely to the trash container .
4
37
10
6
5
30
5
41
4
If the vertical lower first arm with the counterweight (controlled by a properly placed stopper post for limited movement) and the movable verticle link plate affixed at one end joined at the other end to the upper second arm in a horizonal position (controlled by another properlty placed stopper post for limited movement) and stays perpendicular, the horizonal upper second arm , due to its orientation, will remain in a locked position, unable to lift or release said container lid until disengaged as illustrated in FIG..
FIG. 3
FIG. 3
2
2
1
2
3
7
8
9
4
5
3
With reference to , the present disclosure may relate to other automatic gravity inforced lid locking devices for wheeled trash containers . Such devices may be constructed from a durable material such as heavy duty metal or may be constructed from components that remain simular to conventional polyproplene wheeled trash containers. Regardless to what sufficiently durable materials are used, all may open said containers descrbed without interference during the lock and unlock trasitions to accomplish their goals, though it should be stated, that the present disclosure is not mounted directly to a receptacle using drilled holes, welded or press-fitted for attachment like most and if made from durable PVC material, being lighter in weight, it could serve as an inexpensive alternative for manufacturing. Illustrated are the embodiments of the present disclosure's left side to comprise the housing structure in its entirety showing the hooks and , the adjustment chain , both rigid upper and lower pivoting arms and mounted to the right side of the housing structure .
FIG. 4
FIG. 7
1
6
4
5
3
5
4
37
3
Referring now to , an automatic gravity inforced lid locking device is represented in an unlocked position showing both link plate and connected pivotable arms and mounted on the housing structure . The upper horizonal arm and the lower vertical arm with counterweight pivotally moved opposite therefrom the housing structure in a substantially fixed unlocked position by the influence of gravity tilted foward as illustrated in .
FIG. 5
1
4
5
6
4
5
24
27
25
28
26
29
4
5
4
3
16
18
19
17
4
3
6
5
3
20
22
23
21
5
3
6
37
4
38
4
16
3
5
39
3
7
8
1
33
35
34
36
2
42
10
11
13
9
12
3
3
15
14
4
5
30
31
32
3
illustrates an exploded perspective view of the automatic gravity lid locking device as follows: A rigid first pivoting arm is situated in a vertical position affixed to a rigid second pivoting arm situated in a horizontal position by a link plate secured near the end of each said pivoting arm and by means of a bolt and , washer and and nylon nut and in such a fashion that both pivoting arms and can rotate freely. The first pivoting arm is fastened to the housing structure by means of a bolt , washer and nylon nut wherein a spacer nut is situated between said pivoting arm and the housing structure in order to provide proper spacing for the link plate to move freely. The second pivoting arm is fastened to the housing structure also by means of a bolt , washer and nylon nut wherein a spacer nut is situated between said pivoting arm and the housing structure once again to provide proper spacing for the link plate to move freely. A weighted rod is inserted into the bottom of the first pivoting arm and remains in place by means of a cap affixed to the bottom of said pivoting arm and its attachment bolt to the housing structure . The second pivoting arm also comprises a cap mainly for aesthetical purposes. The housing structure also comprises two hooks and in the front of the device affixed by a bolt and and nylon nut and for means of hooking onto a trash container's top edge or lip . A stopper post bolt comprising a bottom washer , a spacer nut , an adjustment chain and a top washer is inserted through the left side of the housing structure pertruding out of the right side of the housing structure comprising a spacer nut and a nylon nut now forming the proper length and means of preventing the first pivoting arm from swinging into an incorrect position. For the same purpose, another such means is provided for the second pivoting arm by including a screw with a protective casing followed by a washer and then affixed to the housing structure on the top right side as illustrated.
1
2
41
2
FIG. 6
Another embodiment of the invention demostrates preferred intended use wherein it is shown in to be hooked and strapped to a a conventional, wheeled, residential type trash container used for garbage or recyclable waste constructed of a plastic type material with a shape that is somewhat rectangular on the upper part and is somewhat cylindrical on the lower part having a hinged lid and is generally indicated by reference numeral .
1
5
41
2
4
37
1
42
2
7
8
41
2
4
6
5
41
In accordance with the invention , said horizonal upper arm having the function of holding down the hingable, closed lid of a trash container is shown in a locked position with a lower vertical arm comprising a counterweight , so that when said device has been interposed between the outer edge wall of said container by hooks and beneath the closed container lid and the receptacle is in an normal upright position on the ground, it then can be tilted forward by a automated garbage truck for dumping. This causes influencing forces of gravity to move the lower vertical arm outwardly thereof, pushing against the link plate in a cooperative manner, disengaging the upper horizonal arm that then recesses upwardly from its locked position holding down the container lid giving way to clear a path of travel to purge the waste.
FIG. 7
2
4
37
6
5
30
41
2
2
41
1
As illustrated in with respect to the trash container being lifted up and tilted forward for dumping, the first arm with the counterweight can pivot outwardly therefrom due to the influence of gravity and push the connecting link plate forward against the second arm forcing it into an unlocked position until it is restricted by a second stopper post so as not to interfere with the hinged lid of the receptacle to achieve clearence for dumping the trash container waste. Thus, when the trash container is emptied and placed back to a substantially upright position and the hinged lid closes, the device returns to a locked position automatically.
Whereas this invention is here illustrated and described with reference to the embodiments thereof presently contemplated as the best mode of carrying out such invention in actual practice, it is to be understood that various changes may be made in adapting the invention to different embodiments without departing from the broader inventive concepts disclosed herein and comprehended by the claims that follow. From the foregoing it will be apparent to those skilled in the art that a variety of changes and modifications can be made in size, shape, type, number and arrangement of parts described without departing from the spirit of this invention and the scope of the attachment claims. | |
Many organizations and teams mistakenly justify context switching as both parallelization and talent resource optimization. Although their goal is efficiency, unknowingly these “hidden priorities” result in not only loss of efficiency (accomplishing things quickly), but also loss of effectiveness (accomplishing the most valuable things).
Strategic & Tactical Prioritization
Prioritization is at the heart of agility. We make tough business prioritization decisions every day in order to deliver what’s most valuable to our customer. The Manifesto for Agile Software Development itself is a set of four values in priority format (“we value X over Y”). The final phrase in the manifesto is especially useful in understanding prioritization: “That is, while there is value in the items on the right, we value the items on the left more.” Prioritization is valuing a particular thing over something else and deciding in what order to address them. Unless your business has infinite bandwidth and resources, it must choose priorities. Not every feature can be implemented today, or even tomorrow. Only one thing at a time can be the next most valuable and deserving of your time, talent and resources.
This is true not only of high-level strategic initiatives but also of the daily tactical task level. Often at the tactical level, if we forget to keep a strategic perspective (i.e. forest for the trees), we end up choosing to do something sub-optimal at the expense of what would have been more optimal. These decisions are based on what we might call “hidden” priorities—hidden somewhere in our subconscious or in habit or in our fears. Bringing these hidden priorities to light can help us determine if they are desirable and, if they aren’t, what actions we can take to overcome them.
This article will examine the hidden priorities in one such decision: context switching vs. parallelization.
Context Switching is NOT Parallelization
At the tactical or task level, context switching is when you start on an item of work before finishing another item of work that is already in process. Agile teams call this practice “thrashing” and avoid it due to the costs of delay it incurs. Parallelization at the task level is multiple team members working on one strategic item of work—they may each be working on different tasks, but collectively are working towards one end-to-end piece of user functionality from the product backlog (many agile teams call these “user stories”). For agile teams, this is often referred to as “swarming,” and they do it to optimize flow and delivery of completed, shippable functionality.
Many organizations mistakenly justify context switching as both parallelization and talent resource optimization. Both involve a one-to-many relationship involving people and items of work but have significantly different outcomes.
Computers (& Humans) Swarm Better Than They Thrash
Computers give us an easy analogy to understand the difference. With a single core machine (back in the day), our PCs would essentially run only one program, one command at a time. So, if you were to start up a program with a long calculation, your mouse would stop responding as the calculation finishes. To remain responsive to the user, the computer frequently paused the existing program to respond to the mouse. To pause a program and start another requires unloading and loading memory, which takes time, and doesn’t directly benefit either program. This is context switching.
In 2001, IBM released the first processor with multiple cores. This allowed the computer to truly run multiple programs at once, since they could each be given different cores. Even better, this allowed one program to split up its work to run on multiple cores, speeding up the work of the program. This is task parallelization. This is one of the main reasons we can do so much with computers today.
What humans can learn from the machines they’ve created is having one worker working concurrently on multiple items of work incurs significant costs. With scrum teams, we have found that the velocity of a sprint without context switching is at least 25% more than a sprint with context switching. Broader scientific and industry research shows that the negative effects of context switching are often even greater:
- Context switching can cost up to 40% of someone’s productive time (American Psychological Association)
- Errors increase in complex work involving a sequence of steps (National Center for Biotechnology Information, U.S. National Library of Medicine)
- Stress is increased (University of California, Irvine)
Gerald Weinberg quantifies the cost of task switching in his book ”Quality Software Management: System Thinking” as follows:
|No. Simultaneous Projects||Percent of Time on Project||Loss to Context Switching|
|1||100%||0%|
|2||40%||20%|
|3||20%||40%|
|4||10%||60%|
|5||5%||75%|
Siloing
Another devious side effect of context switching is “siloing”—when one person takes an entire piece of work on their own, they keep the knowledge independently. That same person will often end up taking all other items that are in that domain moving forward. Down the road, you’ll hit the day where that person is sick or on vacation (no, they don’t even need to get hit by a bus), and the blocking issue in that area can’t be resolved. In addition, team members learn less from each other, aren’t as versatile, and the focus on the sprint goal becomes diluted as people focus on “their own items” and even gold-plating.
Task parallelization, or swarming, generally improves the rate of delivery of a scrum team, as well as their ability to take advantage of the collective and collaborative intelligence of the whole team to adapt when things don’t go according to plan.
Exposing and Removing Hidden Priorities
So, what are the hidden priorities that cause teams and businesses to inappropriately justify context switching? Let’s look at two.
Responsiveness Can Be an Illusion
The first reason for context switching is the same reason computers did it: to give an illusion of responsiveness. Someone wants to be able to tell a customer, “We’ll get right on that.” Unfortunately, ‘getting right on’ something only implies it will get started. It says nothing of finishing it. In other words, in the name of satisfying the customer early, getting work into the system is valued more than getting value out of the system, the latter being what the customer actually wants. This is what most teams actually want as well, but full utilization is that hidden priority that causes them to behave differently.
With the team taking on more work than is sustainable, this overloaded system increases variance in outcomes each sprint and limits their ability to estimate and forecast in the long term. This can lead to hasty and unreasonable commitments, which in turn leads to more schedule overruns and unmet expectations.
Context switching also leads to an increase in defects. With team members having to split their attention between multiple backlog items, nuances and corner cases can fall through the cracks. The team’s definition of done (i.e. what it means to be shippable at the end of a sprint) can get overlooked, cutting corners to make up for lost time. These quality issues turn into more customer requests (sometimes urgent ones) that distracts from innovation and exacerbates the pressure on the team and business expectations.
To combat this, it requires a diligent product owner to be able to work with customers to find the appropriate place in the backlog for new items (or even shuffling older ones), instead of reverting to the “we’ll get right on that” mentality. A good product owner knows how to manage customer expectations so that the team can be allowed to finish the work they started. The scrum master should also watch out for when context switching is occurring and make the issue visible.
Full Utilization Is A Myth
The second reason context switching is used is pressure to show full team member utilization. Utilization is prioritizing individual team member effectiveness over value produced by the team as a whole. Many poor decisions are justified in the name of “100% utilization,” which is a very misleading metric. It doesn’t measure product delivered to the user. It doesn’t measure quality. It doesn’t measure anything that your customer is paying you for. What utilization measures is how many hours a week someone is working, regardless of whether it’s the most valuable work they could be doing. If 100% utilization is your goal, it means that you can justify any behavior as long as people are working 40 hours (often more, unfortunately) a week. You can’t increase delivered product or quality by measuring utilization.
Pursuing utilization often sounds like, “All our remaining work in the sprint is testing, so this developer should just get started on the next sprint.” Or “Bob is a GUI guy, so we’ll have him work ahead on something else that has GUI in it.” The sprint backlog should be the next highest priority items from the product backlog. So, by definition, when team members are working on something outside of the current user story, they are now working on something that isn’t the next highest priority for the business. Utilization devalues business and customer priority. Ironic, isn’t it?
Swarming Isn’t Easy
Teams that are new to swarming may not believe it can be practically done. Common objections include:
- “I don’t know how to do that.” Or, “That’s Megan’s responsibility.” — This can be reflective of an unwillingness to learn, or to do something that someone considers beneath them. In today’s world, learning new skills every day is how we keep up with changing technology and marketplaces. Good teams are constantly learning. The focus should be on contributing to the team’s current goals.
- “I don’t have time to teach John about this.” — This may be trying to establish job security by withholding information. Team members need to understand that job security is really about making the team indispensable, not the individual.
- “There would be too many cooks in the kitchen.” — While sometimes true, the number of people who can work on a user story is often more than originally thought.
Let’s examine how a scrum development team (in the software industry, in this example) might effectively swarm:
- The team elaborates the user story together, which might include the initial design exercise
- One person writes the code
- A second can be writing the unit tests—that’s a separate file, so no stepping on toes there
- A third can be writing integration or system tests —that’s three people already working on one user story
- A fourth can be writing automated user acceptance tests
- Two classes to write? You can add two more people to write code and unit tests
- Documentation updates (wiki, training, user, etc)
- Consider pair programming to front-load peer review for quality
- Writing database upgrade scripts
- Maybe you can think of others…
Development teams are between 3 and 9 people, ideally somewhere in the middle of that range. There is plenty for everyone to do on a single user story.
Yes, this requires collaboration daily during development. That is why scrum is called scrum—a rugby metaphor for team collaboration. Yes, this may mean that sometimes you have a developer working on something that isn’t their strongest area. But what you get is a focus on delivering the highest priority work, earlier and more often, instead of making sure that a worker can show they are busy 100% of their day doing something. And the collaboration ensures that you don’t get siloing, but a team of versatile developers.
The best teams have T-shaped members—individuals with depth or expertise (not the same as specialization) in at least one skill or knowledge area, but have some breadth in multiple areas.
T-shaped developers have versatility and are eager to learn, and the goal is the sprint goal, not utilization. T-shaped members are constantly increasing their ability to contribute to delivering value to the customer.
Escaping “The Thrash”
Consider a few metrics that may be useful to help you avoid the hidden priorities behind context switching. These, as well as others can encourage behaviors that contribute to value delivered instead of rationalizing utilization and individual siloed contributions. Each of the following are outlined in more detail in our “Be Careful What You Measure…” blog:
- Return on investment (ROI): Is desired and intended value actually realized when you deliver new functionality? How early and often is new value delivered to generate a return on your investment?
- Defects in production: If defects are landing in production, do you know why? Is thrashing the reason you’re not satisfying your definition of done to actually mean “shippable”?
- Sprint goal success rate: Are you delivering busy work, getting lots of work started, but not finished? Are you working towards customer goals or on disparate tasks?
- Lead and cycle time: Are you starting more than finishing?
- Skill versatility: Are developers continually increasing their ability to contribute to the sprint goal?
These can all help keep the focus on delivering work, not starting it. In choosing between context switching and swarming, it is important to remember that context switching is about getting work started, swarming is all about getting work finished—delivering value early and often.
By and large, customers are focused on work finished. They don’t care about a feature that they can’t use. They don’t care about utilization. And, they don’t care about how quick you can start. Your customer cares about how quick you can finish and deliver quality. As Steve Jobs purportedly said, “real artists ship.”
If you want help to uncover your hidden priorities to be as effective as you can be, contact an adviser today. | https://platinumedge.com/uncover-the-hidden-priorities-that-make-you-ineffective |
Deriving Quantum Chemical Hamiltonians from Data
- Photophysics of disordered organic systems.
- Deriving quantum chemical hamiltonians from data.
- Design of dyes for imaging of biological systems.
- Design of a light-driven molecular motor.
Overview
Quantum chemistry provides accurate methods for computing the electronic structure of small to medium size molecules, but due to rapid increases in computational cost with system size, applications to large systems is prohibitively expensive. The following two aspects of chemical systems may make computations on large systems computationally feasible:
- Nearsightedness: Interactions between regions of a molecule may be progressively coarse-grained with increasing separation. This enables linear scaling approaches.
- Molecular similarity: Atoms, functional groups, etc. behave similarly in different environments, an assumption which underlies the atom- or group-specific parameters of molecular mechanics and semi-empirical quantum chemistry.
Our goal is to use machine learning algorithms to take better advantage of molecular similarity. Such algorithms can determine useful descriptors of the electronic structures (feature extraction algorithms) and help write the energy and forces as functions of these descriptors (predictors).
The general approach is to first generate a large set of ab initio data on the system of interest, such as a functional group attached to many different molecules or a reaction center in many different environments. This data is then mined for a low-cost model that can reproduce the data with chemical accuracy, but at substantially reduced cost.
In past work, we have applied this approach to develop functional-group specific approaches to electron correlation. The results are analogous to functional-group specific density functionals. This work is summarized in the following video presentation:
Recent work
Figure: Schematic representation of a model that is trained to produce the output of a high level model using information generated from low level models.
More recently, we have developed reliable low-cost quantum mechanical models for use in quantum mechanical/molecular mechanical (QM/MM) simulations of chemical reactions. The H + HF → H2 + F collinear reaction was used as a test case. The approach first generates detailed quantum chemical data for the reaction center in geometries and electrostatic environments that span those expected to arise during the molecular dynamics simulations. For each geometry and environment, both high level (HL) and low level (LL) ab initio calculations are performed. A model is then developed to predict the HL results using only inputs generated from the LL theory. The inputs used here are based on principal component analysis of the LL distributed multipoles (DMs), and the model is a simple linear regression. The DMs are monopoles, dipoles and quadrupoles at each atomic center, and summarize the electronic distribution in a manner that is comparable across basis set. The error in the model is dominated by extrapolation from small to large basis sets, with extrapolation from uncorrelated to correlated methods contributing much less error. A single regression can be used to make predictions for a range of reaction-center geometries and environments. For the trial collinear reaction, separate regressions were developed for the entrance channel, transition region, and exit channel. These models can predict the results of QCISD/6-31++G** computations from HF/3-21G DMs, with an average error for the reaction energy profile of 0.47 kcal/mol. | http://www.chem.cmu.edu/groups/yaron/projects/semimethod.html |
# Interpolation
In the mathematical field of numerical analysis, interpolation is a type of estimation, a method of constructing (finding) new data points based on the range of a discrete set of known data points.
In engineering and science, one often has a number of data points, obtained by sampling or experimentation, which represent the values of a function for a limited number of values of the independent variable. It is often required to interpolate; that is, estimate the value of that function for an intermediate value of the independent variable.
A closely related problem is the approximation of a complicated function by a simple function. Suppose the formula for some given function is known, but too complicated to evaluate efficiently. A few data points from the original function can be interpolated to produce a simpler function which is still fairly close to the original. The resulting gain in simplicity may outweigh the loss from interpolation error and give better performance in calculation process.
## Example
This table gives some values of an unknown function f ( x ) {\displaystyle f(x)} .
Interpolation provides a means of estimating the function at intermediate points, such as x = 2.5. {\displaystyle x=2.5.}
We describe some methods of interpolation, differing in such properties as: accuracy, cost, number of data points needed, and smoothness of the resulting interpolant function.
### Piecewise constant interpolation
The simplest interpolation method is to locate the nearest data value, and assign the same value. In simple problems, this method is unlikely to be used, as linear interpolation (see below) is almost as easy, but in higher-dimensional multivariate interpolation, this could be a favourable choice for its speed and simplicity.
### Linear interpolation
One of the simplest methods is linear interpolation (sometimes known as lerp). Consider the above example of estimating f(2.5). Since 2.5 is midway between 2 and 3, it is reasonable to take f(2.5) midway between f(2) = 0.9093 and f(3) = 0.1411, which yields 0.5252.
Generally, linear interpolation takes two data points, say (xa,ya) and (xb,yb), and the interpolant is given by:
This previous equation states that the slope of the new line between ( x a , y a ) {\displaystyle (x_{a},y_{a})} and ( x , y ) {\displaystyle (x,y)} is the same as the slope of the line between ( x a , y a ) {\displaystyle (x_{a},y_{a})} and ( x b , y b ) {\displaystyle (x_{b},y_{b})}
Linear interpolation is quick and easy, but it is not very precise. Another disadvantage is that the interpolant is not differentiable at the point xk.
The following error estimate shows that linear interpolation is not very precise. Denote the function which we want to interpolate by g, and suppose that x lies between xa and xb and that g is twice continuously differentiable. Then the linear interpolation error is
In words, the error is proportional to the square of the distance between the data points. The error in some other methods, including polynomial interpolation and spline interpolation (described below), is proportional to higher powers of the distance between the data points. These methods also produce smoother interpolants.
### Polynomial interpolation
Polynomial interpolation is a generalization of linear interpolation. Note that the linear interpolant is a linear function. We now replace this interpolant with a polynomial of higher degree.
Consider again the problem given above. The following sixth degree polynomial goes through all the seven points:
Substituting x = 2.5, we find that f(2.5) = ~0.59678.
Generally, if we have n data points, there is exactly one polynomial of degree at most n−1 going through all the data points. The interpolation error is proportional to the distance between the data points to the power n. Furthermore, the interpolant is a polynomial and thus infinitely differentiable. So, we see that polynomial interpolation overcomes most of the problems of linear interpolation.
However, polynomial interpolation also has some disadvantages. Calculating the interpolating polynomial is computationally expensive (see computational complexity) compared to linear interpolation. Furthermore, polynomial interpolation may exhibit oscillatory artifacts, especially at the end points (see Runge's phenomenon).
Polynomial interpolation can estimate local maxima and minima that are outside the range of the samples, unlike linear interpolation. For example, the interpolant above has a local maximum at x ≈ 1.566, f(x) ≈ 1.003 and a local minimum at x ≈ 4.708, f(x) ≈ −1.003. However, these maxima and minima may exceed the theoretical range of the function; for example, a function that is always positive may have an interpolant with negative values, and whose inverse therefore contains false vertical asymptotes.
More generally, the shape of the resulting curve, especially for very high or low values of the independent variable, may be contrary to commonsense; that is, to what is known about the experimental system which has generated the data points. These disadvantages can be reduced by using spline interpolation or restricting attention to Chebyshev polynomials.
### Spline interpolation
Remember that linear interpolation uses a linear function for each of intervals . Spline interpolation uses low-degree polynomials in each of the intervals, and chooses the polynomial pieces such that they fit smoothly together. The resulting function is called a spline.
For instance, the natural cubic spline is piecewise cubic and twice continuously differentiable. Furthermore, its second derivative is zero at the end points. The natural cubic spline interpolating the points in the table above is given by
In this case we get f(2.5) = 0.5972.
Like polynomial interpolation, spline interpolation incurs a smaller error than linear interpolation, while the interpolant is smoother and easier to evaluate than the high-degree polynomials used in polynomial interpolation. However, the global nature of the basis functions leads to ill-conditioning. This is completely mitigated by using splines of compact support, such as are implemented in Boost.Math and discussed in Kress.
### Mimetic interpolation
Depending on the underlying discretisation of fields, different interpolants may be required. In contrast to other interpolation methods, which estimate functions on target points, mimetic interpolation evaluates the integral of fields on target lines, areas or volumes, depending on the type of field (scalar, vector, pseudo-vector or pseudo-scalar).
A key feature of mimetic interpolation is that vector calculus identities are satisfied, including Stokes' theorem and the divergence theorem. As a result, mimetic interpolation conserves line, area and volume integrals. Conservation of line integrals might be desirable when interpolating the electric field, for instance, since the line integral gives the electric potential difference at the endpoints of the integration path. Mimetic interpolation ensures that the error of estimating the line integral of an electric field is the same as the error obtained by interpolating the potential at the end points of the integration path, regardless of the length of the integration path.
Linear, bilinear and trilinear interpolation are also considered mimetic, even if it is the field values that are conserved (not the integral of the field). Apart from linear interpolation, area weighted interpolation can be considered as one of the first mimetic interpolation method to have been developed.
## Function approximation
Interpolation is a common way to approximate functions. Given a function f : → R {\displaystyle f:\to \mathbb {R} } with a set of points x 1 , x 2 , … , x n ∈ {\displaystyle x_{1},x_{2},\dots ,x_{n}\in } one can form a function s : → R {\displaystyle s:\to \mathbb {R} } such that f ( x i ) = s ( x i ) {\displaystyle f(x_{i})=s(x_{i})} for i = 1 , 2 , … , n {\displaystyle i=1,2,\dots ,n} (that is, that s {\displaystyle s} interpolates f {\displaystyle f} at these points). In general, an interpolant need not be a good approximation, but there are well known and often reasonable conditions where it will. For example, if f ∈ C 4 ( ) {\displaystyle f\in C^{4}()} (four times continuously differentiable) then cubic spline interpolation has an error bound given by ‖ f − s ‖ ∞ ≤ C ‖ f ( 4 ) ‖ ∞ h 4 {\displaystyle \|f-s\|_{\infty }\leq C\|f^{(4)}\|_{\infty }h^{4}} where h max i = 1 , 2 , … , n − 1 | x i + 1 − x i | {\displaystyle h\max _{i=1,2,\dots ,n-1}|x_{i+1}-x_{i}|} and C {\displaystyle C} is a constant.
## Via Gaussian processes
Gaussian process is a powerful non-linear interpolation tool. Many popular interpolation tools are actually equivalent to particular Gaussian processes. Gaussian processes can be used not only for fitting an interpolant that passes exactly through the given data points but also for regression; that is, for fitting a curve through noisy data. In the geostatistics community Gaussian process regression is also known as Kriging.
## Other forms
Other forms of interpolation can be constructed by picking a different class of interpolants. For instance, rational interpolation is interpolation by rational functions using Padé approximant, and trigonometric interpolation is interpolation by trigonometric polynomials using Fourier series. Another possibility is to use wavelets.
The Whittaker–Shannon interpolation formula can be used if the number of data points is infinite or if the function to be interpolated has compact support.
Sometimes, we know not only the value of the function that we want to interpolate, at some points, but also its derivative. This leads to Hermite interpolation problems.
When each data point is itself a function, it can be useful to see the interpolation problem as a partial advection problem between each data point. This idea leads to the displacement interpolation problem used in transportation theory.
## In higher dimensions
Multivariate interpolation is the interpolation of functions of more than one variable. Methods include bilinear interpolation and bicubic interpolation in two dimensions, and trilinear interpolation in three dimensions. They can be applied to gridded or scattered data. Mimetic interpolation generalizes to n {\displaystyle n} dimensional spaces where n > 3 {\displaystyle n>3} .
Nearest neighbor Bilinear Bicubic
## In digital signal processing
In the domain of digital signal processing, the term interpolation refers to the process of converting a sampled digital signal (such as a sampled audio signal) to that of a higher sampling rate (Upsampling) using various digital filtering techniques (for example, convolution with a frequency-limited impulse signal). In this application there is a specific requirement that the harmonic content of the original signal be preserved without creating aliased harmonic content of the original signal above the original Nyquist limit of the signal (that is, above fs/2 of the original signal sample rate). An early and fairly elementary discussion on this subject can be found in Rabiner and Crochiere's book Multirate Digital Signal Processing.
## Related concepts
The term extrapolation is used to find data points outside the range of known data points.
In curve fitting problems, the constraint that the interpolant has to go exactly through the data points is relaxed. It is only required to approach the data points as closely as possible (within some other constraints). This requires parameterizing the potential interpolants and having some way of measuring the error. In the simplest case this leads to least squares approximation.
Approximation theory studies how to find the best approximation to a given function by another function from some predetermined class, and how good this approximation is. This clearly yields a bound on how well the interpolant can approximate the unknown function.
## Generalization
If we consider x {\displaystyle x} as a variable in a topological space, and the function f ( x ) {\displaystyle f(x)} mapping to a Banach space, then the problem is treated as "interpolation of operators". The classical results about interpolation of operators are the Riesz–Thorin theorem and the Marcinkiewicz theorem. There are also many other subsequent results. | https://en.wikipedia.org/wiki/Interpolate |
Current Corona Protection Ordinance
The Corona Protection Ordinance (CoronaSchVO)is currently valid from January 16th, 2022, which allows museums to operate. The following rules apply additionally:
- Please note that currently for visiting the Deutsches Bergbau-Museum Bochum the 2G rule for adults and young people aged 16 and over and the 3G rule for children and young people under 16 apply.
- When you visit our museum, we are obliged to checkyourproof.
- For adults and teenagers 16 years and older, this is done via a vaccination card (analog or digital) or a certificate of past illness. Adults who cannot be vaccinated for medical reasons will be allowed access via the "3G" rule upon presentation of a doctor's note and a negative test result. Unvaccinated pregnant women will need their maternity record and a negative test. All test results must be no more than 24 hours old and must have been issued by an approved testing station. PCR tests from recognised laboratories are valid for max. 48 hours. In addition, you must present either a valid photo ID or driver's license for matching.
- Proof for school-age children up to 16 years is provided by school testing or a negative test resultduring the vacations or proof of immunization or recovery.
Children up to school age do not require proof.
- Medical masks are compulsory throughout the museum.
Please note that due to the special local conditions in the Deutsches Bergbau-Museum Bochum, we may deviate from other interpretations of the regulation for the protection of all of us.
- The four tours of the permanent exhibition, the headframe and the visitors’ mine are open to visitors
- Guided tours, educational offerings and events take place. Here, too, the 2G rule for adults and young people aged 16 and over and the 3G rule for children and young people under 16 apply. Depending on the location, the number of participants may be limited.
- During the week, the visitor's mine is only open to visitors on guided tours, at weekends without a guided tour.
- Due to the special ventilation situation in the visitor's mine, we can only allow a limited number of visitors to enter the mine per hour. There is therefore a limit to the number of visitors on all weekdays.
At the moment we ask you to register for the guided tours during the week by telephone on +49 234 5877-220.
At weekends and for the "Triff den Bergmann" (Meet the Miner) format, visitors are taken to the visitor's mine in time slots. It is not possible to register or reserve in advance, but you can do so on site at the ticket office.
- Note on accessibility in the visitor's mine: Access via the lift is only possible at weekends. There may be short waiting times. During the week, the guided tours are still not barrier-free. We apologise for any inconvenience this may cause.
- The museum restaurant "Kumpels" is currently closed due to the pandemic.
- The general rules of the retail trade apply to the museum shop.
You can find an overview of further hygiene measures during your visit here.
We look forward to your visit!
CURRENT INFORMATION
Due to current circumstances, program changes or cancellations may also occur at short notice. Please inform yourself on our homepage and during opening hours by calling the visitor service at +49 234 5877-126 (Tues. to Sun. between 09:30 and 17:30).
INFORMATION FOR VISITORS WITH LIMITED MOBILITY
In view of the current situation, visitors with limited mobility and visitors with prams or buggies are asked to call visitor services. We will hold the door open for you, and activate the lift for your use: +49 234 5877-126. | https://www.bergbaumuseum.de/en/news-detail/current-corona-protection-ordinance |
The flowers of Canada Goldenrod are important nectar and pollen sources for insects, including bees, flies, wasps, and butterflies. Canada Goldenrod is incapable of pollinating itself, a task usually carried out by visiting insects. The plants and their seeds provide food for finches and other birds, and foraging animals (e.g. sheep, cattle, deer, horses). They have also been used by Indigenous people for medicinal and other purposes.
The solitary, hairy stems of this perennial grow 30 to 214 cm tall. Its lance-shaped leaves have finely-toothed margins and hairy undersides. Small flower heads occur along the upper side of each branch, and are arranged into loose, elongated clusters that bloom from the bottom up. Each head consists of several yellow ray and disc florets. The single-seeded fruits have a bristly top to aid in wind-dispersal.
Seeds and/or plants are typically available from greenhouses and seed supply companies specializing in native plants. Canada Goldenrod can be started by seeds, seedlings, or rhizomes. Plants can become weedy, so care should be taken to control spreading. | http://www.prairiepollination.ca/plante-plant/verge_dor_du_canada-canada_goldenrod/ |
Blueberry Muffins...low sugar with egg free, dairy free, gluten free and low gluten options.
Author: Stephen for Natasha
Recipe type: snack, egg free, dairy free, low gluten, low sugar
Prep time:
Cook time:
Total time:
Serves: 10 - 12
Well, we've been testing some combinations for you so we can bring you a muffin with options. You can try them egg free, dairy free, gluten free or low gluten....your choice. This is a low sugar recipe but if you don't find it sweet enough, just add an additional 20 - 30g of sugar. If you don't have a Thermomix just combine the wet ingredients and the dry ingredients separately and then mix together. We hope you enjoy them as much as we do.
Ingredients
- 100g butter or coconut oil (or oil of your choice)
- 1 egg (or 1 tablespoon chia seed soaked in 3 tablespoons of water for 10 mins)
- 180g milk (buttermilk or dairy free milk of your choice)
- 80g brown sugar (we've also used rapadura and coconut palm sugar)
- 1 tablespoon honey (optional)
- 250g self raising flour (or flour of your choice. We've used spelt and wholemeal spelt with the addition of 4 teaspoons of baking powder)
- 120g frozen blueberries
Instructions
- Preheat oven to 200 degrees celcius.
- Place butter or coconut oil in TM bowl and heat for 2 mins at 60 degrees (if you're using an oil in a liquid state you can leave out this step).
- Add milk, egg (or soaked chia) and honey and mix for 5 seconds on speed 3.
- Add flour (and baking powder if you're using plain flour or spelt) and sugar and mix for 10 seconds on speed 4.
- Add blueberries and mix on reverse for 5 seconds on speed 3.
- Spoon (or use an ice cream scoop to transfer) mixture into a greased muffin tin (lined with paper cases if you prefer).
- Bake for approximately 15 mins until golden and springy and a skewer tests clean or with light crumbs.
Notes
The old fashioned ice cream scoop (we have two sizes) is one of our favourite utensils in the the kitchen. We use it to transfer muffin mixture and biscuit dough to baking tins and trays. It's an excellent way to ensure some consistency in size and it just makes the whole job easier. | http://agagfaf.com/2012/05/31/blueberry-muffins-with-egg-free-dairy-free-and-low-gluten-options/ |
Samantha Hossack explains why the US and other NATO members will not follow Hungarian efforts to support the creation of a Kurdish state, giving independence to the Kurds in Iraq?
Sam Hossack
NATO’s Response to Russia: How Energy Concerns Trump Sovereignty
Samantha Hossack examines recent NATO responses to Russia’s actions in Ukraine and what they mean for the Alliance.
Canada’s Contribution to NATO
Sam Hossack examines Canada’s contribution —financial, diplomatic and military— to NATO in light of recent criticisms.
Canadian Foreign Policy Leading the Way for NATO?
Samantha Hossack contemplates the ramifications of recent decisions Canadian foreign policy in regards to its participation in NATO ops.
A World of Difference: The Potential Failures of Geneva II
Samantha Hossack examines the continued struggle for peace in Syria
The Throne View of Canada’s Arctic: How the Throne Speech Identifies Canada’s Arctic Priorities
Samantha Hossack on what the Throne Speech means for Canada’s Arctic plans
Drone-ing On: The Question of Humanitarianism (Part Two)
Samantha Hossack’s second piece about the humanitarian concerns affiliated with drone warfare.
Drone-ing On: The Theoretical Benefits to Drone Warfare
Samantha Hossack discusses the advantages of drone warfare
Canada’s Arctic Priorities: The Best Way to Assure Sovereignty?
Samantha Hossack evaluates Canada’s Arctic Priorities and their impact on Canada’s claim to the region
Russia’s International Image: A Lost Opportunity?
With the recent G20 and upcoming Sochi Olympic Games, Russia has the opportunity to redeem it’s international image. Samantha Hossack’s take on the possibilities therein. | https://natoassociation.ca/category/ncc-authors/sam-hossack/ |
The years building up to an Olympics are always fascinating ones. Calculations are being done across the board to discover how many medals one country can achieve, while for some athletes, the only thing on their mind is qualification for Tokyo.
And despite wall-to-wall sports coverage over the last few months, there has been plenty of activity across the Olympic qualifier spectrum, with Natalya Coyle adding her name to the Irish team and a place in her third consecutive Olympics at the weekend after accomplishing her mission at the European Modern Pentathlon Championships.
It may still be early days but with the team starting to take shape people are probably going to start assessing the chances of an Irish athlete coming home from Tokyo with some extra baggage.
In 2018 alone, 77 Irish athletes secured medals on the world stage. Forty-four of those athletes were female athletes, and seven more were in a mixed group. The rest went towards male athletes.
At the European Championships, hurdler Thomas Barr was the only competitor to come home with a medal having secured bronze in the 400m hurdles, but 2018 was a year for the young guns to prove their worth.
Hungrier
At the Under-20 World Championships in Finland, the quartet of Molly Scott, Gina Akpe-Moses, Ciara Neville and Patience Jumbo-Gula ran an Irish junior record of 43.90 seconds, just .08 seconds off the gold-winning German team who were fastest in qualifying. They are only the third Irish athletes to medal at the U-20 championships, with Ciara Mageean winning silver in the 1500m in 2010, while back in 1994, Antoin Burke won silver in the high jump.
Throw into the mix the even younger and hungrier Irish athletes in Sarah Healy, Sophie O’Sullivan and Rhasidat Adeleke. In 2018, Healy came home with two gold medals from the U-18 European Championships in 1,500m and 3,000m, while O’ Sullivan earned silver in 800m and Adeleke earned gold in 200m. Donegal high jumper Sommer Lecky is also a hope for Ireland too.
Admittedly, it will take some time, nurturing and an ability to maintain these levels to see the senior medals come in, let alone Olympic medals. While the younger athletes have a bit to go in terms of development, race experience and everything needed to be a well-rounded Olympic athlete, those in the athletics world have every right to get excited. For the first time watching the Nationals this year, Irish fans saw a glimmer of hope that won’t die, and many are confident this young crop are here to stay.
Ireland have always competed well in terms of boxing, and 2018 was massive for the sport across all ages, with a total of 29 athletes bringing home medals. The most celebrated include Kellie Harrington at lightweight and Michaela Walsh in the featherweight division. Can either of these boxers replicate the success of Katie Taylor and bring Olympic gold back to Irish soil?
Strong contenders
As for left of field sports that we all simultaneously enjoy every Olympic cycle, sailing and rowing are strong contenders too. Although Annalise Murphy may enjoy poking fun at her unforgettable tears in 2012, she more than made up for it in 2016, adding her name to female Olympic medal holders.
With a new boat and enough sailing around the world completed, Annalise is definitely one to watch along with the rowing contingent, which includes Sanita Puspure (some only know her as dominant) and a bunch of O’Donovans for good measure. Also in the water is backstroke sensation Shane Ryan, who won bronze at the Worlds (25m) and European Championships (50m). As for gymnastics, it is pretty fun to speculate over Rhys McClenaghan and his lovely, lovely pommel horse.
And who can forget about our hockey team, working hard under new coach Sean Dancer, who are looking to storm through the European Championships alongside the likes of Belarus, England and Germany?
The year to date is showing that upward trend in medal hauls with athletes at this year’s European Games athletes coming home with seven medals, six in boxing and one in badminton for the Magee siblings.
Coyle putting her name on the ticket to Tokyo makes this Irish team, hopefully, more excited and well-rounded, with medal hopes coming from a variety of sport as well as boxing and athletics. The Irish team is now full of quality and quantity, and hopefully, more names will be etched in the Irish Olympics history books. | https://www.irishtimes.com/sport/other-sports/natalya-coyle-s-qualification-shows-irish-olympic-team-shaping-up-nicely-1.3986366 |
Protein disulfide isomerase (PDI) is an endoplasmic reticulum (ER)-resident oxidoreductase chaperone that catalyzes the maturation of disulfide-bond-containing proteins. S-nitrosylated PDI (SNO-PDI), which is the chemical modification of its catalytic cysteine in response to nitrosative stress, has been found in post-mortem brains of Parkinson's and Alzheimer's disease victims along with synphilin-1:alpha-synuclein protein aggregates (called Lewy Bodies). Additional studies in cells revealed that levels of SNO-PDI formation directly co-related with the aggregation and accumulation of the minor but critical Parkinsonian biomarker synphilin-1 in a NO-sensitive manner. While SNO-PDI formation leads to the accumulation of polyubiquitinated proteins, expression of native PDI (non-SNO-PDI) attenuates these effects in a Parkinsonian cell model. These data show that PDI is neuroprotective and underscore the need for functional preservation of PDI's catalytic activity as a key preventative approach to pathogenesis of nitrosative-stress-related Parkinson's. The data also suggest the involvement of PDI dysfunction in the pathogenesis of neuropathies such as the Lewy Body Variant of Alzheimer's (LBVAD) and Alzheimer's. However, there is still a gap in studies designed to determine whether other neurotoxicity-related biomarkers accumulate as a function of SNO-PDI formation. We hypothesize that SNO-PDI formation may provoke aggregation of alpha-synuclein, the major Parkinsonian and LBVAD biomarker protein and Lewy-body constituent. As a corollary to our hypothesis, it is possible that strategies designed to prevent SNO-PDI formation are neuroprotective and prophylactic to Parkinson's. The hypothesis will be tested by examining the aggregation of alpha-synuclein as a function of nitrosative insult in a cell line model. Furthermore, the translational feasibility of ellagic acid, Na-betahydroxybutyrate and Ferrostatin-1 analogs, which our lab has preliminarily demonstrated as being neuroprotective, will be assayed in a rotenone rat Parkinson model. The overall objective of this project it to lay the foundation for long-term work involving the development of pharmacologically relevant small molecule therapies that are neuroprotective by mitigating the effects of oxidative and nitrosative stress.
| |
How do you get myositis?
Myositis means inflammation of the muscles that you use to move your body. An injury, infection, or autoimmune disease can cause it.
Who is most likely to get myositis?
Adults between the ages of 30 and 60, and children between the ages of 5 and 15 are more likely to get myositis.
Can inflammatory myopathy be prevented?
Myopathy Prevention
There are no known ways to prevent myopathy. Muscle inflammation may be caused by an allergic reaction, exposure to a toxic substance or medicine, another disease such as cancer or rheumatic conditions, or a virus or other infectious agent.
Can myositis come and go?
The onset of symptoms usually occurs gradually over a period of months. Occasionally, however, symptoms can develop rapidly over a period of days. Symptoms may also come and go for no apparent reason. The main symptom associated with polymyositis is muscle weakness.
What does myositis feel like?
Myositis is the name for a group of rare conditions. The main symptoms are weak, painful or aching muscles. This usually gets worse, slowly over time. You may also trip or fall a lot, and be very tired after walking or standing.
How long can you live with myositis?
More than 95 percent of those with DM, PM, and NM are still alive more than five years after diagnosis. Many experience only one period of acute illness in their lifetime; others struggle with symptoms for years. One of the biggest problems in treating myositis is obtaining an accurate diagnosis.
How do I know if I have myositis?
Myositis usually begins gradually, but can take a variety of forms. Sometimes the first sign is an unusual rash. Sometimes patients may start to trip or fall more frequently. Other signs include muscle weakness and pain, intense fatigue, and trouble climbing stairs or reaching over the head.
When does myositis start?
The risk of developing IBM increases with age and usually appears in patients over the age of 50; however, patients may develop symptoms as early as their 30s.
Does exercise help myositis?
Physical exercise has been shown to reduce inflammation, reduce fatigue, increase stamina, and build muscle, even in patients with myositis. Indeed, exercise is currently the only treatment recommendation for patients with inclusion body myositis.
What autoimmune disease causes leg swelling?
Myositis (my-o-SY-tis) is a rare type of autoimmune disease that inflames and weakens muscle fibers. Autoimmune diseases occur when the body’s own immune system attacks itself. In the case of myositis, the immune system attacks healthy muscle tissue, which results in inflammation, swelling, pain, and eventual weakness.
What autoimmune disease causes myositis?
There are four types of autoimmune myositis:
- Polymyositis.
- Dermatomyositis.
- Necrotizing immune-mediated myopathies.
- Inclusion body myositis.
How do they test for myositis?
Muscle and skin biopsy are often the most definitive way to diagnose myositis diseases. Small samples of muscle tissue show abnormalities in muscles, including inflammation, damage, and abnormal proteins. For those with skin symptoms, doctors often biopsy a bit of skin to examine for characteristic abnormalities. | https://computersurgery.org/orthopedics/can-myositis-be-prevented.html |
400 U.S. 74 (1970)
DUTTON, WARDEN
v.
EVANS.
No. 10.
Supreme Court of United States.
Argued October 15, 1969.
Reargued October 15, 1970
Decided December 15, 1970
APPEAL FROM THE UNITED STATES COURT OF APPEALS FOR THE FIFTH CIRCUIT.
*75 Alfred L. Evans, Jr., Assistant Attorney General of Georgia, reargued the cause for appellant. With him on the brief were Arthur K. Bolton, Attorney General, and Marion O. Gordon and Mathew Robins, Assistant Attorneys General.
Robert B. Thompson reargued the cause and filed a brief for appellee.
Solicitor General Griswold, by invitation of the Court, argued the cause for the United States as amicus curiae on the reargument. With him on the brief were Assistant Attorney General Wilson, Jerome M. Feit, Beatrice Rosenberg, and Roger A. Pauley.
*76 MR. JUSTICE STEWART announced the judgment of the Court and an opinion in which THE CHIEF JUSTICE, MR. JUSTICE WHITE, and MR. JUSTICE BLACKMUN join.
Early on an April morning in 1964, three police officers were brutally murdered in Gwinnett County, Georgia. Their bodies were found a few hours later, handcuffed together in a pine thicket, each with multiple gunshot wounds in the back of the head. After many months of investigation, Georgia authorities charged the appellee, Evans, and two other men, Wade Truett and Venson Williams, with the officers' murders. Evans and Williams were indicted by a grand jury; Truett was granted immunity from prosecution in return for his testimony.
Evans pleaded not guilty and exercised his right under Georgia law to be tried separately. After a jury trial, he was convicted of murder and sentenced to death.[1] The judgment of conviction was affirmed by the Supreme Court of Georgia,[2] and this Court denied certiorari.[3] Evans then brought the present habeas corpus proceeding in a federal district court, alleging, among other things, that he had been denied the constitutional right of confrontation at his trial. The District Court denied the writ,[4] but the Court of Appeals for the Fifth Circuit reversed, holding that Georgia had, indeed, denied Evans the right, guaranteed by the Sixth and Fourteenth Amendments, "to be confronted by the witnesses against him."[5] From that judgment an appeal was brought to this Court, and we noted probable jurisdiction.[6] The *77 case was originally argued last Term, but was set for reargument. 397 U. S. 1060.
In order to understand the context of the constitutional question before us, a brief review of the proceedings at Evans' trial is necessary. The principal prosecution witness at the trial was Truett, the alleged accomplice who had been granted immunity. Truett described at length and in detail the circumstances surrounding the murder of the police officers. He testified that he, along with Evans and Williams, had been engaged in switching the license plates on a stolen car parked on a back road in Gwinnett County when they were accosted by the three police officers. As the youngest of the officers leaned in front of Evans to inspect the ignition switch on the car, Evans grabbed the officer's gun from its holster. Evans and Williams then disarmed the other officers at gunpoint, and handcuffed the three of them together. They then took the officers into the woods and killed them by firing several bullets into their bodies at extremely close range. In addition to Truett, 19 other witnesses testified for the prosecution.[7] Defense counsel was given full opportunity to cross-examine each witness, and he exercised that opportunity with respect to most of them.
One of the 20 prosecution witnesses was a man named Shaw. He testified that he and Williams had been fellow prisoners in the federal penitentiary in Atlanta, Georgia, at the time Williams was brought to Gwinnett County to be arraigned on the charges of murdering the police officers. Shaw said that when Williams was returned to the penitentiary from the arraignment, he had asked Williams: "How did you make out in court?" and that Williams had responded, "If it hadn't been for that dirty son-of-a-bitch Alex Evans, we wouldn't be in this now." Defense counsel objected to the introduction *78 of this testimony upon the ground that it was hearsay and thus violative of Evans' right of confrontation. After the objection was overruled, counsel cross-examined Shaw at length.
The testimony of Shaw relating what he said Williams had told him was admitted by the Georgia trial court, and its admission upheld by the Georgia Supreme Court, upon the basis of a Georgia statute that provides: "After the fact of conspiracy shall be proved, the declarations by any one of the conspirators during the pendency of the criminal project shall be admissible against all."[8] As the appellate court put it:
" `The rule is that so long as the conspiracy to conceal the fact that a crime has been committed or the identity of the perpetrators of the offense continues, the parties to such conspiracy are to be considered so much a unit that the declarations of either are admissible against the other.' The defendant, and his co-conspirator, Williams, at the time this statement was made, were still concealing their identity, keeping secret the fact that they had killed the deceased, if they had, and denying their guilt. There was evidence sufficient to establish a prima facie case of conspiracy to steal the automobile and the killing of the deceased by the conspirators while carrying out the conspiracy, and the statement by Williams made after the actual commission of the crime, but while the conspiracy continued was admissible."[9] (Citations omitted.)
This holding was in accord with a consistent line of Georgia decisions construing the state statute. See, e. g., Chatterton v. State, 221 Ga. 424, 144 S. E. 2d 726, *79 cert. denied, 384 U. S. 1015; Burns v. State, 191 Ga. 60, 73, 11 S. E. 2d 350, 358.
It was the admission of this testimony of the witness Shaw that formed the basis for the appellee's claim in the present habeas corpus proceeding that he had been denied the constitutional right of confrontation in the Georgia trial court. In upholding that claim, the Court of Appeals for the Fifth Circuit regarded its duty to be "not only to interpret the framers' original concept in light of historical developments, but also to translate into due-process terms the constitutional boundaries of the hearsay rule."[10] (Footnotes omitted.) The court upheld the appellee's constitutional claim because it could find no "salient and cogent reasons" for the exception to the hearsay rule Georgia applied in the present case, an exception that the court pointed out was broader than that applicable to conspiracy trials in the federal courts.[11]
The question before us, then, is whether in the circumstances of this case the Court of Appeals was correct in holding that Evans' murder conviction had to be set aside because of the admission of Shaw's testimony. In considering this question, we start by recognizing that this Court has squarely held that "the Sixth Amendment's right of an accused to confront the witnesses against him is . . . a fundamental right . . . made obligatory on the States by the Fourteenth Amendment." Pointer v. Texas, 380 U. S. 400, 403. See also Douglas v. Alabama, 380 U. S. 415; Brookhart v. Janis, 384 U. S. 1; Barber v. Page, 390 U. S. 719; Roberts v. Russell, 392 U. S. 293; Illinois v. Allen, 397 U. S. 337; California v. Green 399 U. S. 149. But that is no more than the beginning of our inquiry.
*80 I
It is not argued, nor could it be, that the constitutional right to confrontation requires that no hearsay evidence can ever be introduced. In the Pointer case itself, we referred to the decisions of this Court that have approved the admission of hearsay:
"This Court has recognized the admissibility against an accused of dying declarations, Mattox v. United States, 146 U. S. 140, 151, and of testimony of a deceased witness who has testified at a former trial, Mattox v. United States, 156 U. S. 237, 240-244. See also Dowdell v. United States, supra, 221 U. S., at 330; Kirby v. United States, supra, 174 U. S., at 61. . . . There are other analogous situations which might not fall within the scope of the constitutional rule requiring confrontation of witnesses."[12]
The argument seems to be, rather, that in any given case the Constitution requires a reappraisal of every exception to the hearsay rule, no matter how long established, in order to determine whether, in the words of the Court of Appeals, it is supported by "salient and cogent reasons." The logic of that position would seem to require a constitutional reassessment of every established hearsay exception, federal or state, but in the present case it is argued only that the hearsay exception applied by Georgia is constitutionally invalid because it does not identically conform to the hearsay exception applicable to conspiracy trials in the federal courts. Appellee does not challenge and we do not question the validity of the coconspirator exception applied in the federal courts.
*81 That the two evidentiary rules are not identical must be readily conceded. It is settled that in federal conspiracy trials the hearsay exception that allows evidence of an out-of-court statement of one conspirator to be admitted against his fellow conspirators applies only if the statement was made in the course of and in furtherance of the conspiracy, and not during a subsequent period when the conspirators were engaged in nothing more than concealment of the criminal enterprise. Lutwak v. United States, 344 U. S. 604; Krulewitch v. United States, 336 U. S. 440. The hearsay exception that Georgia applied in the present case, on the other hand, permits the introduction of evidence of such an out-of-court statement even though made during the concealment phase of the conspiracy.
But it does not follow that because the federal courts have declined to extend the hearsay exception to include out-of-court statements made during the concealment phase of a conspiracy, such an extension automatically violates the Confrontation Clause. Last Term in California v. Green, 399 U. S. 149, we said:
"Our task in this case is not to decide which of these positions, purely as a matter of the law of evidence, is the sounder. The issue before us is the considerably narrower one of whether a defendant's constitutional right `to be confronted with the witnesses against him' is necessarily inconsistent with a State's decision to change its hearsay rules . . . . While it may readily be conceded that hearsay rules and the Confrontation Clause are generally designed to protect similar values, it is quite a different thing to suggest that the overlap is complete and that the Confrontation Clause is nothing more or less than a codification of the rules of hearsay and their exceptions as they existed historically at common law. Our decisions have never established *82 such a congruence; indeed, we have more than once found a violation of confrontation values even though the statements in issue were admitted under an arguably recognized hearsay exception. The converse is equally true: merely because evidence is admitted in violation of a long-established hearsay rule does not lead to the automatic conclusion that confrontation rights have been denied." Id., at 155-156 (citations and footnote omitted).
These observations have particular force in the present case. For this Court has never indicated that the limited contours of the hearsay exception in federal conspiracy trials are required by the Sixth Amendment's Confrontation Clause. To the contrary, the limits of this hearsay exception have simply been defined by the Court in the exercise of its rule-making power in the area of the federal law of evidence.[13] It is clear that the limited scope of the hearsay exception in federal conspiracy trials is a product, not of the Sixth Amendment, but of the Court's "disfavor" of "attempts to broaden the already pervasive and wide-sweeping nets of conspiracy prosecutions." Grunewald v. United States, 353 U. S. 391, 404. As Grunewald, Krulewitch, and other cases in this Court make clear, the evidentiary rule is intertwined, not only with the federal substantive law of conspiracy, but also with such related issues as the impact of the statute of limitations upon conspiracy prosecutions.
*83 In the case before us such policy questions are not present. Evans was not prosecuted for conspiracy in the Georgia court, but for the substantive offense of murder.[14] At his trial the State permitted the introduction of evidence under a long-established and well-recognized rule of state law.[15] We cannot say that the evidentiary rule applied by Georgia violates the Constitution merely because it does not exactly coincide with the hearsay exception applicable in the decidedly different context of a federal prosecution for the substantive offense of conspiracy.
II
It is argued, alternatively, that in any event Evans' conviction must be set aside under the impact of our recent decisions that have reversed state court convictions because of the denial of the constitutional right of confrontation. The cases upon which the appellee Evans primarily relies are Pointer v. Texas, supra; Douglas *84 v. Alabama, supra; Brookhart v. Janis, supra; Barber v. Page, supra; and Roberts v. Russell, supra.
In the Pointer case it appeared that a man named Phillips had been the victim of a robbery in Texas. At a preliminary hearing, Phillips "as chief witness for the State gave his version of the alleged robbery in detail, identifying petitioner as the man who had robbed him at gunpoint." 380 U. S., at 401. Pointer had no lawyer at this hearing and did not try to cross-examine Phillips. At Pointer's subsequent trial the prosecution was permitted to introduce the transcript of Phillips' testimony given at the preliminary hearing. Thus, as this Court held, the State's "use of the transcript of that statement at the trial denied petitioner any opportunity to have the benefit of counsel's cross-examination of the principal witness against him." 380 U. S., at 403. The Douglas case, decided the same day as Pointer, involved an even more flagrant violation of the defendant's right of confrontation. For at Douglas' trial the prosecutor himself was permitted to read an "entire document" purporting to be an accomplice's written confession after the accomplice had refused to testify in reliance upon his privilege against compulsory self-incrimination. "The statements from the document as read by the Solicitor recited in considerable detail the circumstances leading to and surrounding the alleged crime; of crucial importance, they named the petitioner as the person who fired the shotgun blast which wounded the victim." 380 U. S., at 417. In reversing Douglas' conviction, this Court pointed out that the accomplice's reliance upon the privilege against compulsory self-incrimination "created a situation in which the jury might improperly infer both that the statement had been made and that it was true." 380 U. S., at 419. Yet, since the prosecutor was "not a witness, the inference from his reading that [the accomplice] made the statement could not be *85 tested by cross-examination. Similarly, [the accomplice] could not be cross-examined on a statement imputed to but not admitted by him." Ibid.
Brookhart v. Janis and Barber v. Page are even further afield. In Brookhart it appeared that the petitioner had been "denied the right to cross-examine at all any witnesses who testified against him," and that, additionally, "there was introduced as evidence against him an alleged confession, made out of court by one of his co-defendants. . . who did not testify in court." 384 U. S., at 4. The only issue in the case was one of waiver, since the State properly conceded that such a wholesale and complete "denial of cross-examination without waiver . . . would be constitutional error of the first magnitude . . . ." 384 U. S., at 3. In Barber the "principal evidence" against the petitioner was a transcript of preliminary hearing testimony admitted by the trial judge under an exception to the hearsay rule that, by its terms, was applicable only if the witness was "unavailable." This hearsay exception "has been explained as arising from necessity . . . ." 390 U. S., at 722, and we decided only that Oklahoma could not invoke that concept to use the preliminary hearing transcript in that case without showing "a good-faith effort" to obtain the witness' presence at the trial. Id., at 725.
In Roberts v. Russell we held that the doctrine of Bruton v. United States, 391 U. S. 123, was applicable to the States and was to be given retroactive effect. But Bruton was a case far different from the one now before us. In that case there was a joint trial of the petitioner and a codefendant, coincidentally named Evans, upon a charge of armed postal robbery. A postal inspector testified that Evans had confessed to him that Evans and the petitioner had committed the robbery. This evidence was, concededly, wholly inadmissible against the petitioner. Evans did not testify. Although the trial judge *86 instructed the jury to disregard the evidence of Evans' confession in considering the question of the petitioner's guilt, we reversed the petitioner's conviction. The primary focus of the Court's opinion in Bruton was upon the issue of whether the jury in the circumstances presented could reasonably be expected to have followed the trial judge's instructions. The Court found that "[t]he risk of prejudice in petitioner's case was even more serious than in Douglas," because "the powerfully incriminating extrajudicial statements of a codefendant, who stands accused side-by-side with the defendant, are deliberately spread before the jury in a joint trial." 391 U. S., at 127, 135-136. Accordingly, we held that "in the context of a joint trial we cannot accept limiting instructions as an adequate substitute for petitioner's constitutional right of cross-examination." 391 U. S., at 137. There was not before us in Bruton "any recognized exception to the hearsay rule," and the Court was careful to emphasize that "we intimate no view whatever that such exceptions necessarily raise questions under the Confrontation Clause." 391 U. S., at 128 n. 3.
It seems apparent that the Sixth Amendment's Confrontation Clause and the evidentiary hearsay rule stem from the same roots.[16] But this Court has never equated the two,[17] and we decline to do so now. We confine ourselves, instead, to deciding the case before us.
*87 This case does not involve evidence in any sense "crucial" or "devastating," as did all the cases just discussed. It does not involve the use, or misuse, of a confession made in the coercive atmosphere of official interrogation, as did Douglas, Brookhart, Bruton, and Roberts. It does not involve any suggestion of prosecutorial misconduct or even negligence, as did Pointer, Douglas, and Barber. It does not involve the use by the prosecution of a paper transcript, as did Pointer, Brookhart, and Barber. It does not involve a joint trial, as did Bruton and Roberts. And it certainly does not involve the wholesale denial of cross-examination, as did Brookhart.
In the trial of this case no less than 20 witnesses appeared and testified for the prosecution. Evans' counsel was given full opportunity to cross-examine every one of them. The most important witness, by far, was the eyewitness who described all the details of the triple murder and who was cross-examined at great length. Of the 19 other witnesses, the testimony of but a single one is at issue here. That one witness testified to a brief conversation about Evans he had with a fellow prisoner in the Atlanta Penitentiary. The witness was vigorously and effectively cross-examined by defense counsel.[18] His testimony, which was of peripheral significance at most, was admitted in evidence under a co-conspirator exception to the hearsay rule long established under state statutory law. The Georgia statute can *88 obviously have many applications consistent with the Confrontation Clause, and we conclude that its application in the circumstances of this case did not violate the Constitution.
Evans was not deprived of any right of confrontation on the issue of whether Williams actually made the statement related by Shaw. Neither a hearsay nor a confrontation question would arise had Shaw's testimony been used to prove merely that the statement had been made. The hearsay rule does not prevent a witness from testifying as to what he has heard; it is rather a restriction on the proof of fact through extrajudicial statements. From the viewpoint of the Confrontation Clause, a witness under oath, subject to cross-examination, and whose demeanor can be observed by the trier of fact, is a reliable informant not only as to what he has seen but also as to what he has heard.[19]
The confrontation issue arises because the jury was being invited to infer that Williams had implicitly identified Evans as the perpetrator of the murder when he blamed Evans for his predicament. But we conclude that there was no denial of the right of confrontation as to this question of identity. First, the statement contained no express assertion about past fact, and consequently it carried on its face a warning to the jury against giving the statement undue weight. Second, Williams' personal knowledge of the identity and role of the other participants in the triple murder is abundantly established by Truett's testimony and by Williams' prior conviction. It is inconceivable that cross-examination could have shown that Williams was not in a position to know *89 whether or not Evans was involved in the murder. Third, the possibility that Williams' statement was founded on faulty recollection is remote in the extreme. Fourth, the circumstances under which Williams made the statement were such as to give reason to suppose that Williams did not misrepresent Evans' involvement in the crime. These circumstances go beyond a showing that Williams had no apparent reason to lie to Shaw. His statement was spontaneous, and it was against his penal interest to make it. These are indicia of reliability which have been widely viewed as determinative of whether a statement may be placed before the jury though there is no confrontation of the declarant.
The decisions of this Court make it clear that the mission of the Confrontation Clause is to advance a practical concern for the accuracy of the truth-determining process in criminal trials by assuring that "the trier of fact [has] a satisfactory basis for evaluating the truth of the prior statement." California v. Green, 399 U. S., at 161. Evans exercised, and exercised effectively, his right to confrontation on the factual question whether Shaw had actually heard Williams make the statement Shaw related. And the possibility that cross-examination of Williams could conceivably have shown the jury that the statement, though made, might have been unreliable was wholly unreal.
Almost 40 years ago, in Snyder v. Massachusetts, 291 U. S. 97, Mr. Justice Cardozo wrote an opinion for this Court refusing to set aside a state criminal conviction because of the claimed denial of the right of confrontation. The closing words of that opinion are worth repeating here:
"There is danger that the criminal law will be brought into contemptthat discredit will even touch the great immunities assured by the Fourteenth Amendmentif gossamer possibilities of prejudice *90 to a defendant are to nullify a sentence pronounced by a court of competent jurisdiction in obedience to local law, and set the guilty free." 291 U. S., at 122.
The judgment of the Court of Appeals is reversed, and the case is remanded to that court for consideration of the other issues presented in this habeas corpus proceeding.[20]
It is so ordered.
MR. JUSTICE BLACKMUN, whom THE CHIEF JUSTICE joins, concurring.
I join MR. JUSTICE STEWART'S opinion. For me, however, there is an additional reason for the result.
The single sentence attributed in testimony by Shaw to Williams about Evans, and which has prolonged this litigation, was, in my view and in the light of the entire record, harmless error if it was error at all. Furthermore, the claimed circumstances of its utterance are so incredible that the testimony must have hurt, rather than helped, the prosecution's case. On this ground alone, I could be persuaded to reverse and remand.
Shaw testified that Williams made the remark at issue when Shaw "went to his room in the hospital" and asked Williams how he made out at a court hearing on the preceding day. On cross-examination, Shaw stated that he was then in custody at the federal penitentiary in Atlanta; that he worked as a clerk in the prison hospital; that Williams was lying on the bed in his *91 room and facing the wall; that he, Shaw, was in the hall and not in the room when he spoke with Williams; that the door to the room "was closed"; that he spoke through an opening about 10 inches square; that the opening "has a piece of plate glass, window glass, just ordinary window glass, and a piece of steel mesh"; that this does not impede talking through the door; and that one talks in a normal voice when he talks through that door. Shaw conceded that when he had testified at Williams' earlier trial, he made no reference to the glass in the opening in the door.
Carmen David Mabry, called by the State, testified that he was with the United States Public Health Service and stationed at the Atlanta Penitentiary. He described the opening in the door to Williams' room and said that it contained a glass "and over that is a wire mesh, heavy steel mesh"; that he has "never tried to talk through the door"; that, to his knowledge, he has never heard "other people talking through the door"; that, during his 11 years at the hospital, the glass has not been out of the door; and that the hospital records disclosed that it had not been out.
I am at a loss to understand how any normal jury, as we must assume this one to have been, could be led to believe, let alone be influenced by, this astonishing account by Shaw of his conversation with Williams in a normal voice through a closed hospital room door. I note, also, the Fifth Circuit's description of Shaw's testimony as "somewhat incredible" and as possessing "basic incredibility." 400 F. 2d, at 828 n. 4.
In saying all this, I am fully aware that the Fifth Circuit panel went on to observe, in the footnote just cited, "[W]e are convinced that it cannot be called harmless." And Justice Quillian, in sole dissent on the direct appeal to the Supreme Court of Georgia, stated, "[I]t obviously was prejudicial to the defendant." 222 Ga. *92 392, 408; 150 S. E. 2d 240, 251. However, neither the Georgia Superior Court judge who tried the case nor the Federal District Judge who held the hearing on Evans' petition for federal habeas concluded that prejudicial error was present. Also, we do not know the attitude of the Georgia Supreme Court majority, for they decided the issue strictly upon the pronounced limits of the long-established Georgia hearsay rule, 222 Ga., at 402; 150 S. E. 2d, at 248, and presumably had no occasion to touch upon any alternative ground such as harmlessness. I usually would refrain from passing upon an issue of this kind adversely to a federal court of appeals, but when the trial judges do not rule, I would suppose that we are as free to draw upon the cold record as is the appellate court.
I add an observation about corroboration. Marion Calvin Perry, another federal prisoner and one who admitted numerous past convictions, including "larceny of automobiles," testified without objection that he had known Williams and Evans for about 10 years, and Truett for about two years; that he spoke with Williams and Evans some 25 or 30 days prior to the murders of the three police officers; that Williams owed him money; that he and Williams talked by telephone "[a]bout me stealing some cars for him"; that Williams told him that "Alex [Evans] would know what kind of car he [Williams] would want"; that a few days later "me and Alex talked about cars and I told him I didn't want to mess with Venson [Williams]"; that Evans said, "if I got any, he said I could get them for him"; that seven or eight days before the murders Williams asked him by telephone whether he, Perry, "still had the Oldsmobile switch"; that the week of the murders he argued with Evans about how much he should receive for each stolen car; that six days after the murders he saw Evans at a filling station; that they talked about the murders; that "I said if I wanted to know who did it, I would see *93 mine and your friend"; and that Evans "got mad as hell" and "told me if I thought I knowed anything about it to keep my damn mouth shut."
Another witness, Lawrence H. Hartman, testified that his 1963 red Oldsmobile hardtop was stolen from his home in Atlanta the night of April 16, 1964 (the murders took place on the early morning of April 17). He went on to testify that the 1963 Oldsmobile found burning near the scene of the tragedy was his automobile. There is testimony in the record as to the earlier acquisition by Evans and Williams of another wrecked Oldsmobile of like model and color; as to the towing of that damaged car by a wrecker manned by Williams and Evans; and as to the replacement of good tires on a Chevrolet occupied by Williams, Evans, and Truett, with recapped tires then purchased by them.
This record testimony, it seems to me, bears directly and positively on the Williams-Evans-Truett car-stealing conspiracy and accomplishments and provides indisputable confirmation of Evans' role. The requirements of the Georgia corroboration rule were fully satisfied and Shaw's incredible remark fades into practical and legal insignificance.
The error here, if one exists, is harmless beyond a reasonable doubt. Chapman v. California, 386 U. S. 18, 21-25; Harrington v. California, 395 U. S. 250.
MR. JUSTICE HARLAN, concurring in the result.
Not surprisingly the difficult constitutional issue presented by this case has produced multiple opinions. MR. JUSTICE STEWART finds Shaw's testimony admissible because it is "wholly unreal" to suggest that cross-examination would have weakened the effect of Williams' statement on the jury's mind. MR. JUSTICE BLACKMUN, while concurring in this view, finds admission of the statement to be harmless, seemingly because he deems Shaw's testimony so obviously fabricated that no normal jury *94 would have given it credence. MR. JUSTICE MARSHALL answers both suggestions to my satisfaction, but he then adopts a position that I cannot accept. He apparently would prevent the prosecution from introducing any out-of-court statement of an accomplice unless there is an opportunity for cross-examination, and this regardless of the circumstances in which the statement was made and regardless of whether it is even hearsay.
The difficulty of this case arises from the assumption that the core purpose of the Confrontation Clause of the Sixth Amendment is to prevent overly broad exceptions to the hearsay rule. I believe this assumption to be wrong. Contrary to things as they appeared to me last Term when I wrote in California v. Green, 399 U. S. 149, 172 (1970). I have since become convinced that Wigmore states the correct view when he says:
"The Constitution does not prescribe what kinds of testimonial statements (dying declarations, or the like) shall be given infra-judicially,this depends on the law of Evidence for the time being,but only what mode of procedure shall be followedi. e. a cross-examining procedurein the case of such testimony as is required by the ordinary law of Evidence to be given infra-judicially." 5 J. Wigmore, Evidence § 1397, at 131 (3d ed. 1940) (footnote omitted).
The conversion of a clause intended to regulate trial procedure into a threat to much of the existing law of evidence and to future developments in that field is not an unnatural shift, for the paradigmatic evil the Confrontation Clause was aimed attrial by affidavit[1]can be *95 viewed almost equally well as a gross violation of the rule against hearsay and as the giving of evidence by the affiant out of the presence of the accused and not subject to cross-examination by him. But however natural the shift may be, once made it carries the seeds of great mischief for enlightened development in the law of evidence.
If one were to translate the Confrontation Clause into language in more common use today, it would read: "In all criminal prosecutions, the accused shall enjoy the right to be present and to cross-examine the witnesses against him." Nothing in this language or in its 18th century equivalent would connote a purpose to control the scope of the rules of evidence. The language is particularly ill-chosen if what was intended was a prohibition on the use of any hearsaythe position toward which my Brother MARSHALL is being driven, although he does not quite yet embrace it.
Nor am I now content with the position I took in concurrence in California v. Green, supra, that the Confrontation Clause was designed to establish a preferential rule, requiring the prosecutor to avoid the use of hearsay where it is reasonably possible for him to do soin other words, to produce available witnesses. Further consideration in the light of facts squarely presenting the issue, as Green did not, has led me to conclude that this is not a happy intent to be attributed to the Framers absent compelling linguistic or historical evidence pointing in that direction. It is common ground that the historical understanding of the clause furnishes no solid guide to adjudication.[2]
A rule requiring production of available witnesses would significantly curtail development of the law of *96 evidence to eliminate the necessity for production of declarants where production would be unduly inconvenient and of small utility to a defendant. Examples which come to mind are the Business Records Act, 28 U. S. C. §§ 1732-1733, and the exceptions to the hearsay rule for official statements, learned treatises, and trade reports. See, e. g., Uniform Rules of Evidence 63 (15), 63 (30), 63 (31); Gilstrap v. United States, 389 F. 2d 6 (CA5 1968) (business records); Kay v. United States, 255 F. 2d 476 (CA4 1958) (laboratory analysis). If the hearsay exception involved in a given case is such as to commend itself to reasonable men, production of the declarant is likely to be difficult, unavailing, or pointless. In unusual cases, of which the case at hand may be an example, the Sixth Amendment guarantees federal defendants the right of compulsory process to obtain the presence of witnesses, and in Washington v. Texas, 388 U. S. 14 (1967), this Court held that the Fourteenth Amendment extends the same protection to state defendants.[3]
Regardless of the interpretation one puts on the words of the Confrontation Clause, the clause is simply not well designed for taking into account the numerous factors that must be weighed in passing on the appropriateness of rules of evidence. The failure of MR. JUSTICE STEWART'S opinion to explain the standard by which it tests Shaw's statement, or how this standard can be squared with the seemingly absolute command of the clause, bears witness to the fact that the clause is being set a task for which it is not suited. The task is far more appropriately performed under the aegis of the Fifth and *97 Fourteenth Amendments' commands that federal and state trials, respectively, must be conducted in accordance with due process of law. It is by this standard that I would test federal and state rules of evidence.[4]
It must be recognized that not everything which has been said in this Court's cases is consistent with this position. However, this approach is not necessarily inconsistent with the results that have been reached. Of the major "confrontation" decisions of this Court, seven involved the use of prior-recorded testimony.[5] In the absence of countervailing circumstances, introduction of such evidence would be an affront to the core meaning of the Confrontation Clause. The question in each case, therefore, was whether there had been adequate "confrontation" to satisfy the requirement of the clause. Regardless of the correctness of the results, the holding that the clause was applicable in those situations is consistent with the view of the clause I have taken.
Passing on to the other principal cases, Dowdell v. United States, 221 U. S. 325, 330 (1911), held that the Confrontation Clause did not prohibit the introduction of "[d]ocumentary evidence to establish collateral facts, *98 admissible under the common law." While this was characterized as an exception to the clause, rather than a problem to which the clause did not speak, the result would seem correct. Brookhart v. Janis, 384 U. S. 1 (1966), and Smith v. Illinois, 390 U. S. 129 (1968), involved restrictions on the right to cross-examination or the wholesale denial of that right. Douglas v. Alabama, 380 U. S. 415 (1965), is perhaps most easily dealt with by viewing it as a case of prosecutorial misconduct. Alternatively, I would be prepared to hold as a matter of due process that a confession of an accomplice resulting from formal police interrogation cannot be introduced as evidence of the guilt of an accused, absent some circumstance indicating authorization or adoption. The exclusion of such evidence dates at least from Tong's Case, Kelyng 17, 18-19, 84 Eng. Rep. 1061, 1062 (K. B. 1663), and is universally accepted. This theory would be adequate to account for the results of both Douglas and Bruton v. United States, 391 U. S. 123 (1968).
The remaining confrontation case of significance is Kirby v. United States, 174 U. S. 47 (1899). In that case a record of conviction of three men for theft was introduced at Kirby's trial. The judge instructed the jury that this judgment was prima facie evidence that the goods which Kirby was accused of receiving from the three men were in fact stolen. This Court reversed, holding that since the judgment was the sole evidence of the fact of theft, Kirby had been denied his right of confrontation. In my view this is not a confrontation case at all, but a matter of the substantive law of judgments. Accord, 4 Wigmore, supra, § 1079, at 133. Indeed, the Kirby Court indicated that lack of confrontation was not at the heart of its objection when it said *99 that the record would have been competent evidence of the fact of conviction. The correctness of the result in Kirby can hardly be doubted, but it was, I think, based on the wrong legal theory.
Judging the Georgia statute here challenged by the standards of due process, I conclude that it must be sustained. Accomplishment of the main object of a conspiracy will seldom terminate the community of interest of the conspirators. Declarations against that interest evince some likelihood of trustworthiness. The jury, with the guidance of defense counsel, should be alert to the obvious dangers of crediting such testimony. As a practical matter, unless the out-of-court declaration can be proved by hearsay evidence, the facts it reveals are likely to remain hidden from the jury by the declarant's invocation of the privilege against self-incrimination.[6] In light of such considerations, a person weighing the necessity for hearsay evidence of the type here involved against the danger that a jury will give it undue credit might reasonably conclude that admission of the evidence would increase the likelihood of just determinations of truth. Appellee has not suggested that Shaw's testimony possessed any peculiar characteristic that would lessen the force of these general considerations and require, as a constitutional matter, that the trial judge exercise residual discretion to exclude the evidence as unduly inflammatory. *100 Exclusion of such statements, as is done in the federal courts, commends itself to me, but I cannot say that it is essential to a fair trail. The Due Process Clause requires no more.
On the premises discussed in this opinion, I concur in the reversal of the judgment below.
MR. JUSTICE MARSHALL, whom MR. JUSTICE BLACK, MR. JUSTICE DOUGLAS, and MR. JUSTICE BRENNAN join, dissenting.
Appellee Evans was convicted of first-degree murder after a trial in which a witness named Shaw was allowed to testify, over counsel's strenuous objection, about a statement he claimed was made to him by Williams, an alleged accomplice who had already been convicted in a separate trial.[1] According to Shaw, the statement, which implicated both Williams and Evans in the crime, was made in a prison conversation immediately after Williams' arraignment. Williams did not testify nor was he called as a witness. Nevertheless, the Court today concludes that admission of the extrajudicial statement attributed to an alleged partner in crime did not deny Evans the right "to be confronted with the witnesses against him" guaranteed by the Sixth and Fourteenth Amendments to the Constitution. In so doing, the majority reaches a result completely inconsistent with recent opinions of this Court, especially Douglas v. Alabama, 380 U. S. 415 (1965), and Bruton v. United States, 391 U. S. 123 (1968). In my view, those cases fully apply here and establish a clear violation of Evans' constitutional rights.
*101 In Pointer v. Texas, 380 U. S. 400 (1965), this Court first held that "the Sixth Amendment's right of an accused to confront the witnesses against him is . . . a fundamental right and is made obligatory on the States by the Fourteenth Amendment." Id., at 403. That decision held constitutionally inadmissible a statement offered against a defendant at a state trial where the statement was originally made at a preliminary hearing under circumstances not affording the defendant an adequate opportunity for cross-examination. Indeed, we have since held that even cross-examination at a prior hearing does not satisfy the confrontation requirement, at least where the witness who made the statement is available to be called at trial. Barber v. Page, 390 U. S. 719 (1968). "The right to confrontation is basically a trial right. It includes both the opportunity to crossexamine and the occasion for the jury to weigh the demeanor of the witness." Id., at 725.
In Douglas v. Alabama, supra, this Court applied the principles of Pointer to a case strikingly similar to this one. There, as here, the State charged two defendants with a crime and tried them in separate trials. There, as here, the State first prosecuted one defendant (Loyd) and then used a statement by him in the trial of the second defendant (Douglas). Although the State called Loyd as a witness, an appeal from his conviction was pending and he refused to testify on the ground that doing so would violate his Fifth Amendment privilege against self-incrimination.
Without reaching the question whether the privilege was properly invoked,[2] the Court held that the prosecutor's *102 reading of Loyd's statement in a purported attempt to refresh his memory denied Douglas' right to confrontation. "Loyd could not be cross-examined on a statement imputed to but not admitted by him." 380 U. S., at 419. Of course, Douglas was provided the opportunity to cross-examine the officers who testified regarding Loyd's statement. "But since their evidence tended to show only that Loyd made the confession, cross-examination of them . . . could not substitute for cross-examination of Loyd to test the truth of the statement itself."[3]Id., at 420. Surely, the same reasoning compels the exclusion of Shaw's testimony here. Indeed, the only significant difference between Douglas and this case, insofar as the denial of the opportunity to cross-examine is concerned, is that here the State did not even attempt to call Williams to testify in Evans' trial. He was plainly available to the State, and for all we know he would have willingly testified, at least with regard to his alleged conversation with Shaw.[4]
Finally, we have applied the reasoning of Douglas to hold that, "despite instructions to the jury to disregard *103 the implicating statements in determining the codefendant's guilt or innocence, admission at a joint trial of a defendant's extrajudicial confession implicating a codefendant violated the codefendant's right of cross-examination secured by the Confrontation Clause of the Sixth Amendment." Roberts v. Russell, 392 U. S. 293 (1968), giving retroactive effect in both state and federal trials to Bruton v. United States, 391 U. S. 123 (1968). Thus Williams' alleged statement, an extrajudicial admission made to a fellow prisoner, could not even have been introduced against Williams if he had been tried in a joint trial with Evans.
The teaching of this line of cases seems clear: Absent the opportunity for cross-examination, testimony about the incriminating and implicating statement allegedly made by Williams was constitutionally inadmissible in the trial of Evans.
MR. JUSTICE STEWART'S opinion for reversal characterizes as "wholly unreal" the possibility that cross-examination of Williams himself would change the picture presented by Shaw's account. A trial lawyer might well doubt, as an article of the skeptical faith of that profession, such a categorical prophecy about the likely results of careful cross-examination. Indeed, the facts of this case clearly demonstrate the necessity for fuller factual development which the corrective test of cross-examination makes possible. The plurality for reversal pigeonholes the out-of-court statement that was admitted in evidence as a "spontaneous" utterance, hence to be believed. As the Court of Appeals concluded, however, there is great doubt that Williams even made the statement attributed to him.[5] Moreover, *104 there remains the further question what, if anything, Williams might have meant by the remark that Shaw recounted. MR. JUSTICE STEWART'S opinion concedes that the remark is ambiguous. Plainly it stands as an accusation of some sort: "If it hadn't been for . . . Evans," said Williams, according to Shaw, "we wouldn't be in this now." At his trial Evans himself gave unsworn testimony to the effect that the murder prosecution might have arisen from enmities that Evans' own law enforcement activities had stirred up in the locality. Did Williams' accusation relate to Evans as a man with powerful and unscrupulous enemies, or Evans as a murderer? MR. JUSTICE STEWART'S opinion opts for the latter interpretation, for it concludes that Williams' remark was "against his penal interest" and hence to be believed. But at this great distance from events, no one can be certain. The point is that absent cross-examination of Williams himself, the jury was left with only the unelucidated, apparently damning, and patently damaging accusation as told by Shaw.
Thus we have a case with all the unanswered questions that the confrontation of witnesses through cross-examination is meant to aid in answering: What did the declarant say, and what did he mean, and was it the truth? If Williams had testified and been cross-examined, Evans' counsel could have fully explored these and other matters. The jury then could have evaluated the statement in the light of Williams' testimony and demeanor. As it was, however, the State was able to use Shaw to present the damaging evidence and thus to avoid confronting Evans with the person who allegedly gave witness against him. I had thought that this was precisely what the Confrontation Clause as applied to the States in Pointer and our other cases prevented.
Although MR. JUSTICE STEWART'S opinion for reversal concludes that there was no violation of Evans' right of *105 confrontation, it does so in the complete absence of authority or reasoning to explain that result. For example, such facts as that Williams' alleged statement was not made during official interrogation, was not in transcript form, and was not introduced in a joint trialthough they differentiate some of the cases are surely irrelevant. Other cases have presented each of these factors,[6] and no reason is offered why the right of confrontation could be so limited.
Nor can it be enough that the statement was admitted in evidence "under a long-established and well-recognized rule of state law." MR. JUSTICE STEWART'S opinion surely does not mean that a defendant's constitutional right of confrontation must give way to a state evidentiary rule. That much is established by our decision in Barber v. Page, supra which held unconstitutional the admission of testimony in accordance with a rule similarly well recognized and long established. However, the plurality for reversal neither succeeds in distinguishing that case nor considers generally that there are inevitably conflicts between Pointer and state evidentiary rules. Rather, it attempts to buttress its conclusion merely by announcing a reluctance to equate evidentiary hearsay rules and the Confrontation Clause.[7]
*106 The Court of Appeals, however, was not of the view that the Confrontation Clause implies unrelenting hostility to whatever evidence may be classified as hearsay. Nor did that court hold that States must conform their evidentiary rules to the hearsay exceptions applicable in federal conspiracy trials. While it did note that this case does not in reality even involve the traditional hearsay rule and its so-called coconspirators exception,[8] that was not the basis for its decision. Rather, the Court of Appeals found in the admission of an incriminatory and inculpating statement attributed to an alleged accomplice who was not made available for cross-examination what it termed an obvious abridgment of Evans' right of confrontation. Since the State presented no satisfactory justification for the denial of confrontation, cf. Pointer v. Texas, 380 U. S., at 407, the Court of Appeals *107 held that under Douglas v. Alabama and this Court's other cases Evans was denied his constitutional rights.
Surely the Constitution requires at least that much when the State denies a defendant the right to confront and cross-examine the witnesses against him in a criminal trial. In any case, that Shaw's testimony was admitted in accordance with an established rule of state law cannot aid my Brethren in reaching their conclusion. Carried to its logical end, justification of a denial of the right of confrontation on that basis would provide for the wholesale avoidance of this Court's decisions in Douglas and Burton,[9] decisions which MR. JUSTICE STEWART'S opinion itself reaffirms. Indeed, if that opinion meant what it says, it would come very close to establishing in reverse the very equation it seeks to avoidan equation that would give any exception to a state hearsay rule a "permanent niche in the Constitution" in the form of an exception to the Confrontation Clause as well.
Finally, the plurality for reversal apparently distinguishes the present case on the ground that it "does not involve evidence in any sense `crucial' or `devastating.' " *108 Despite the characterization of Shaw's testimony as "of peripheral significance at most," however, the possibility of its prejudice to Evans was very real. The outcome of Evans' trial rested, in essence, on whether the jury would believe the testimony of Truett with regard to Evans' role in the murder. Truett spoke as an admitted accomplice who had been immunized from prosecution. Relying on Georgia law, not federal constitutional law, the trial judge instructed the jury that "you cannot lawfully convict upon the testimony of an accomplice alone. . . . [T]he testimony of an accomplice must be corroborated. . . . [T]he corroboration . . . must be such as to connect the defendant with the criminal act." The State presented the testimony of a number of other witnesses, in addition to that of the alleged accomplice that tended to corroborate Evans' guilt. But Shaw's account of what Williams supposedly said to him was undoubtedly a part of that corroborating evidence.[10]
*109 Indeed, MR. JUSTICE STEWART'S opinion does not itself upset the Court of Appeals' finding that the admission of Shaw's testimony, if erroneous, could not be considered harmless. Beyond and apart from the question of harmless error, MR. JUSTICE STEWART undertakes an inquiry, the purpose of which I do not understand, into whether the evidence admitted is "crucial" or "devastating." The view is, apparently, that to require the exclusion of evidence falling short of that high standard of prejudice would bring a moment of clamor against the Bill of Rights. I would eschew such worries and confine the inquiry to the traditional questions: Was the defendant afforded the right to confront the witnesses against him? And, if not, was the denial of his constitutional right harmless beyond a reasonable doubt?
The fact is that Evans may well have been convicted in part by an incriminatory and implicating statement attributed to an alleged accomplice who did not testify and who consequently could not be questioned regarding the truth or meaning of that statement. The Court of Appeals correctly recognized that the Confrontation Clause prohibits such a result, whether the statement is introduced under the guise of refreshing a witness' recollection as in Douglas v. Alabama, against a codefendant with a limiting instruction as in Bruton v. United States, or in accordance with some other evidentiary rule as here.
I am troubled by the fact that the plurality for reversal, unable when all is said to place this case beyond the principled reach of our prior decisions, shifts its ground and begins a hunt for whatever "indicia of reliability" may cling to Williams' remark, as told by Shaw. Whether Williams made a "spontaneous" statement "against his penal interest" is the very question that should have been tested by cross-examination of Williams *110 himself. If "indicia of reliability" are so easy to come by, and prove so much, then it is only reasonable to ask whether the Confrontation Clause has any independent vitality at all in protecting a criminal defendant against the use of extrajudicial statements not subject to cross-examination and not exposed to a jury assessment of the declarant's demeanor at trial.[11] I believe the Confrontation Clause has been sunk if any out-of-court statement bearing an indicium of a probative likelihood can come in, no matter how damaging the statement may be or how great the need for the truth-discovering test of cross-examination. Cf. California v. Green, 399 U. S. 149, 161-162 (1970). Our decisions from Pointer and Douglas to Bruton and Roberts require more than this meager inquiry. Nor is the lame "indicia" approach necessary to avoid a rampaging Confrontation Clause that tramples all flexibility and innovation in a state's law of evidence. That specter is only a specter.[12] To decide this case I need not go beyond hitherto settled Sixth and Fourteenth Amendment law to consider generally what effect, if any, the Confrontation Clause has on the common-law hearsay rule and its exceptions, since no issue of such global dimension is presented. Cf. Bruton v. United States, 391 U. S., at 128 n. 3. The incriminatory extrajudicial statement of an alleged accomplice is so inherently prejudicial that it cannot be introduced unless there is an opportunity to cross-examine the declarant, whether or not his statement *111 falls within a genuine exception to the hearsay rule.
In my view, Evans is entitled to a trial in which he is fully accorded his constitutional guarantee of the right to confront and cross-examine all the witnesses against him. I would affirm the judgment of the Court of Appeals and let this case go back to the Georgia courts to be tried without the use of this out-of-court statement attributed by Shaw to Williams.
NOTES
[1] The parties agree that this death sentence cannot be carried out. See n. 20, infra.
[2] Evans v. State, 222 Ga. 392, 150 S. E. 2d 240.
[3] 385 U. S. 953.
[4] The opinion of the District Court is unreported.
[5] Evans v. Dutton, 400 F. 2d 826, 827.
[6] 393 U. S. 1076. Since, as will appear, the Court of Appeals held that a Georgia statute relied upon by the State at the trial was unconstitutional as applied, there can be no doubt of the right of appeal to this Court. 28 U. S. C. § 1254 (2).
[7] Three of these were rebuttal witnesses. There were four defense witnesses, and Evans himself made a lengthy unsworn statement.
[8] Ga. Code Ann. § 38-306 (1954).
[9] Evans v. State, 222 Ga. 392, 402, 150 S. E. 2d 240, 248.
[10] 400 F. 2d, at 829.
[11] 400 F. 2d, at 830, 831.
[12] Pointer v. Texas, 380 U. S., at 407. See also Salinger v. United States, 272 U. S. 542, 548.
[13] See 18 U. S. C. § 3771. Fed. Rule Crim. Proc. 26 provides:
"In all trials the testimony of witnesses shall be taken orally in open court, unless otherwise provided by an act of Congress or by these rules. The admissibility of evidence and the competency and privileges of witnesses shall be governed, except when an act of Congress or these rules otherwise provide, by the principles of the common law as they may be interpreted by the courts of the United States in the light of reason and experience."
See Hawkins v. United States, 358 U. S. 74.
[14] We are advised that at the time of Evans' trial Georgia did not recognize conspiracy as a separate, substantive criminal offense.
[15] The Georgia rule is hardly unique. See, e. g., Reed v. People, 156 Colo. 450, 402 P. 2d 68; Dailey v. State, 233 Ala. 384,171 So. 729; State v. Roberts, 95 Kan. 280, 147 P. 828. See also 2 F. Wharton, Criminal Evidence § 430 (12th ed. 1955):
"The acts and declarations of a conspirator are admissible against a co-conspirator when they are made during the pendency of the wrongful act, and this includes not only the perpetration of the offense but also its subsequent concealment. . . .
.....
"The theory for the admission of such evidence is that persons who conspire to commit a crime, and who do commit a crime, are as much concerned, after the crime, with their freedom from apprehension, as they were concerned, before the crime, with its commission: the conspiracy to commit the crime devolves after the commission thereof into a conspiracy to avoid arrest and implication."
The existence of such a hearsay exception in the evidence law of many States was recognized in Krulewitch, supra. 336 U. S., at 444.
[16] It has been suggested that the constitutional provision is based on a common-law principle that had its origin in a reaction to abuses at the trial of Sir Walter Raleigh. F. Heller, The Sixth Amendment 104 (1951).
[17] See Note, Confrontation and the Hearsay Rule, 75 Yale L. J. 1434:
"Despite the superficial similarity between the evidentiary rule and the constitutional clause, the Court should not be eager to equate them. Present hearsay law does not merit a permanent niche in the Constitution; indeed, its ripeness for reform is a unifying theme of evidence literature. From Bentham to the authors of the Uniform Rules of Evidence, authorities have agreed that present hearsay law keeps reliable evidence from the courtroom. If Pointer has read into the Constitution a hearsay rule of unknown proportions, reformers must grapple not only with centuries of inertia but with a constitutional prohibition as well." Id., at 1436. (Footnotes omitted.)
[18] This cross-examination was such as to cast serious doubt on Shaw's credibility and, more particularly, on whether the conversation which Shaw related ever took place.
[19] Of course Evans had the right to subpoena witnesses, including Williams, whose testimony might show that the statement had not been made. Counsel for Evans informed us at oral argument that he could have subpoenaed Williams but had concluded that this course would not be in the best interests of his client.
[20] It was conceded at oral argument that the death penalty imposed in this case cannot be carried out, because the jury was qualified under standards violative of Witherspoon v. Illinois, 391 U. S. 510. The Court of Appeals for the Fifth Circuit has already set aside, under Witherspoon, the death sentence imposed upon Venson Williams, Evans' alleged accomplice. See Williams v. Dutton, 400 F. 2d 797, 804-805.
[1] See California v. Green, supra, at 179 (concurring opinion): historically, "the Confrontation Clause was meant to constitutionalize a barrier against flagrant abuses, trial by anonymous accusers, and absentee witnesses."
[2] See id., at 175-179, especially 176 n. 8 (concurring opinion).
[3] Although the fact is not necessary to my conclusion, I note that counsel for Evans conceded at oral argument that he could have secured Williams' presence to testify, but decided against it. Tr. of Oral Arg. 51, 55.
[4] Reliance on the Due Process Clauses would also have the virtue of subjecting rules of evidence to constitutional scrutiny in civil and criminal trials alike. It is exceedingly rare for the common law to make admissibility of evidence turn on whether the proceeding is civil or criminal in nature. See 1 Wigmore, supra, § 4, at 16-17. This feature of our jurisprudence is a further indication that the Confrontation Clause, which applies only to criminal prosecutions, was never intended as a constitutional standard for testing rules of evidence.
[5] Reynolds v. United States, 98 U. S. 145 (1879); Mattox v. United States, 156 U. S. 237 (1895); Motes v. United States 178 U. S. 458 (1900); West v. Louisiana, 194 U. S. 258 (1904); Pointer v. Texas, 380 U. S. 400 (1965); Barber v. Page, 390 U. S. 719 (1968); California v. Green, 399 U. S. 149 (1970).
[6] Quite apart from Malloy v. Hogan, 378 U. S. 1 (1964), Georgia has long recognized the privilege. The Georgia Constitution of 1877, Art. I, § 1, ¶ VI, provided that: "No person shall be compelled to give testimony tending in any manner to criminate himself," and the same language appears in the present state constitution. Ga. Const. of 1945, Art. I, § 1, ¶ VI. The right had previously been recognized as a matter of common law, even in civil trials. See, e. g., Marshall v. Riley, 7 Ga. 367 (1849).
[1] Shaw had been a witness at Williams' trial; his testimony was fully anticipated and was objected to both before and after its admission.
[2] This same questionwhich presents a fundamental conflict between a defendant's Sixth Amendment rights and a witness' Fifth Amendment privilegemight have been present here had the State called Williams to testify. Under a view that would make availability of a declarant the only concern of confrontation, see California v. Green, 399 U. S. 149, 172-189 (1970) (HARLAN, J., concurring), the State's right or duty to compel a codefendant's testimony, by timing of trials and use of testimonial immunity, would seemingly have to be decided. See Comment, Exercise of the Privilege Against Self-Incrimination by Witnesses and Codefendants: The Effect Upon the Accused, 33 U. Chi. L. Rev. 151, 165 (1965).
[3] Cf. Brookhart v. Janis, 384 U. S. 1, 4 (1966).
[4] My Brother STEWART comments that Evans might have brought Williams to the courthouse by subpoena. Defense counsel did not do so, believing that Williams would stand on his right not to incriminate himself. Tr. of Oral Arg. 55. Be that as it may, it remains that the duty to confront a criminal defendant with the witnesses against him falls upon the State, and here the State was allowed to introduce damaging evidence without running the risks of trial confrontation. Cf. n. 2, supra.
[5] After considering Shaw's testimony and other evidence submitted at the trial, the Court of Appeals concluded that Shaw's account of his conversation with Williams was notable for "its basic incredibility." 400 F. 2d 826, 828 n. 4.
[6] For example, Pointer involved only the second, and that one was not present in either Bruton or Roberts.
[7] Constitutionalization of "all common-law hearsay rules and their exceptions," California v. Green, 399 U. S., at 174 (concurring opinion), would seem to be a prospect more frightening than real. Much of the complexity afflicting hearsay rules comes from the definition of hearsay as an out-of-court statement presented for the truth of the matter stateda definition nowhere adopted by this Court for confrontation purposes. Rather, the decisions, while looking to availability of a declarant, Barber v. Page, supra recognize that "cross-examination is included in the right of an accused in a criminal case to confront the witnesses against him," Pointer v. Texas, 380 U. S., at 404, and that admission in the absence of cross-examination of certain types of suspect and highly damaging statements is one of the "threats to a fair trial" against which "the Confrontation Clause was directed," Bruton v. United States, 391 U. S., at 136.
[8] Evans was not charged with conspiracy nor could he have been under Georgia law. The "conspiracy" element came in as part of the State's evidentiary law, part of which goes far beyond the traditional hearsay exception even as it exists with regard to the "concealment phase" in some jurisdictions. Indeed, Williams' alleged statement itself negates the notion that Evans had authorized Williams to speak or had assumed the risk in order to achieve an unlawful aim through concert of effort. It is difficult to conceive how Williams could be part of a conspiracy to conceal the crime when all the alleged participants were in custody and he himself had already been arraigned. As this Court stated in Fiswick v. United States, 329 U. S. 211, 217 (1946), an "admission by one co-conspirator after he has been apprehended is not in any sense a furtherance of the criminal enterprise. It is rather a frustration of it." One lower court in Georgia has adopted essentially this reasoning in reversing a conviction where testimony similar to that objected to in this case was admitted. See Green v. State, 115 Ga. App. 685, 155 S. E. 2d 655 (1967). But see n. 9, infra.
[9] The Georgia rule involved here, which apparently makes admissible all pre-trial statements and admissions of an alleged accomplice or coconspirator, inevitably conflicts with this Court's decisions regarding the Confrontation Clause. See Darden v. State, 172 Ga. 590, 158 S. E. 414 (1931), and Mitchell v. State, 86 Ga. App. 292, 71 S. E. 2d 756 (1952), where confessions of codefendants not on trial were held admissible. Indeed, the Georgia Supreme Court seems to have resolved this conflict in favor of the state rule by erroneously concluding that this Court's decisions are based on the federal hearsay rule concerning "a confession by one of the co-conspirators after he has been apprehended." Pinion v. State, 225 Ga. 36, 37, 165 S. E. 2d 708, 709-710 (1969). See also Park v. State, 225 Ga. 618,170 S. E. 2d 687 (1969), petition for cert. filed, November 4, 1969, No. 57, O. T. 1970 (renumbered).
[10] The trial judge's instructions left no doubt that the statement attributed to Williams could provide the necessary corroboration. See Trial Record 412-413. Indeed, the prejudicial impact of Shaw's testimony is graphically revealed simply by juxtaposing two quotations. First, there is characterization in MR. JUSTICE STEWART'S opinion of Shaw's testimony, a characterization that I find fair albeit studiedly mild: "[T]he jury was being invited to infer that Williams had implicitly identified Evans as the perpetrator of the murder. . . ." (Emphasis added.) Second, there is the trial judge's charge on corroboration of accomplice testimony: "Slight evidence from an extraneous source identifying the accused as a participator in the criminal act will be sufficient corroboration of an accomplice to support a verdict." (Emphasis added.) In the light of the charge and on consideration of the whole record of Evans' trial, it is impossible for me to believe "beyond a reasonable doubt" that the error complained of did not contribute to the verdict obtained. Chapman v. California, 386 U. S. 18, 24 (1967); Harrington v. California, 395 U. S. 250, 251 (1969).
[11] MR. JUSTICE HARLAN answers this question with directness by adopting, to decide this case, his view of due process which apparently makes no distinction between civil and criminal trials, and which would prohibit only irrational or unreasonable evidentiary rulings. Needless to say, I cannot accept the view that Evans' constitutional rights should be measured by a standard concededly having nothing to do with the Confrontation Clause.
[12] See n. 7, supra.
| |
Masters Theses. Ishaan Grover, A semantics based computational model for word learning, 2018. S. M. Media Arts and Sciences, MIT. Nikhita Singh, Talking machines: democratizing the design of voice-based agents for the home, 2018.S. M. Media Arts and Sciences, MIT. Pedro Reynolds-Cuéllar, The role of social robots in fostering human empathy : a cross-cultural exploration, 2018.
Robotics Ph.D. Program | School of Interactive Computing
Aug 26, 2017 · Everyone knows about the Miracle on the Hudson, but you might not know that approximately one in every 2,000 flights experiences a bird strike.A vast majority of these don’t end with a plane ditched in a river or a bird piercing a cockpit window, but this common problem is a constant threat to aviation safety.. In fact, between 1960 and 2014, bird strikes were responsible for the destruction
JKRobots.com
Everyone on our professional essay writing team is an expert in academic research and in APA, MLA, Chicago, Harvard citation formats. Your project arrives fully formatted and ready to submit. The research behind the writing is always 100% original, and the writing is guaranteed Phd Thesis Robot free of …
PhD in Robotics: Information for Doctoral Students
The TU/e is searching for:Two fully-funded PhD positions on impact-aware robot manipulation TU/e, the Netherlands at the Department of Mechanical Engineering, Dynamics and Control.
PhD Thesis. Vision-based gesture recognition in a robot
Cubelets were invented by Modular Robotics CEO Eric Schweikardt as part of his PhD thesis. When a Cubelet wants to communicate with another Cubelet to which it is not direclty connected, it passes data through the intermediate Cubelets. Specs. ROBOTS is a product of IEEE Spectrum, the flagship publication of the IEEE,
Ph.D. Program in Robotics
robots.engin.umich.edu
Theses and Dissertations - Robotics - LibGuides at
Student Services: Bill & Melinda Gates Center, Box 352355 3800 E Stevens Way NE Seattle, WA 98195-2355
Degree Requirements | Michigan Robotics
The AIRO group (Artificial Intelligence and Robotics) of the IDLab at Ghent University is looking for two PhD students to help with our research on social robotics and machine learning. Social robots are robots which interact with people in a natural manner by, for example, using speech, gestures, facial expressions and language. As social robots…
PhD Theses – Intelligent Robots and Systems Group
Apr 10, 2019 · Robotics Thesis/Dissertation Advice. question. Hello all, I am a first year robotics engineering combined MS PhD student (5 year program) and I am having trouble finding a research topic that falls into my interests. My adviser suggested that I apply for grants for rehabilitation robotics, but my interest is in space robotics, although it could
Doctoral Program in Robotics - The Robotics Institute
Here you will find a list of relevant PhD thesis in reverse chronological order by year; with each year the entries are alphabetical by author. When attaching your PhD thesis, please follow the formatting of previous entries.. Please upload your thesis to the wiki. Do not just add a link to your homepage! This page is supposed to serve as a repository that will still exist after your homepage
Two PhD positions in social robotics « Tony Belpaeme
Disclaimer: is the online writing service that offers custom written Phd Thesis Robotics papers, including research papers, thesis papers, essays and others. Online writing service includes the research material as well, but these services are for assistance purposes only. All papers from this agency should be properly referenced.
Phd Thesis Robotics
control strategies for robots in contact a dissertation submitted to the department of aeronautics & astronautics and the committee on graduate studies of stanford university in partial fulfillment of the requirements for the degree of doctor of philosophy jaeheung park march 2006
PhD Theses | Algorithms for Planning and Control of Robot
Course Requirements MS PhD. The Robotics Masters (MS) degree program requires completion of 30 credits of letter-graded coursework including directed study for 3 to 6 credits. PhD programs have very similar course requirements. PhD students earn a Masters degree as part of their PhD program.
RobotSpaceBrain - A Discovery of Science & Art
PhD in Robotics Engineering. A Ph.D. program in robotics engineering can advance your skills in designing and developing robots for numerous applications. Check the prerequisites for enrolling in this program, and review the coursework. Get career outlook and salary info for robotics engineers.
Mobile Robot Phd Thesis
PhD Theses. 2018. Distributed State Estimation and Control of Autonomous Quadrotor Formations Using Exclusively Onboard Resources. Duarte Dias. 2018. An Institutional Economics-based Approach to Distributed Robotic Systems. José Nuno Pereira. 2014. 2013. An Integrated Bayesian Approach to Multi-Robot Cooperative Perception. Aamir Ahmad.
Ph.D. Candidates for 2019-2020 | Computer Science
In addition to performing original research culminating in a doctoral thesis, students complete 36 hours of coursework and focus on three of five core robotics areas: Mechanics, Controls, Perception, Human-Robot Interaction (HRI), and Artificial Intelligence (AI) & Autonomy.
Thesis – Soft Robotics Research
The inter-disciplinary robotics program offers Master of Science, and Doctor of Philosophy degrees in Robotics. Master's degree candidates may pursue thesis or nonthesis options. The PhD program prepares students for careers in industry, research laboratories or universities.
Phd positions on impact-aware robot manipulation
Jun 09, 2017 · Ph.D. Candidate Defends Thesis on Coordinating Dynamics in Human-Robot Teams. Jun 9, 2017. Tariq Iqbal with robot. It's a long way from his native Bangladesh to La Jolla, California, where Ph.D. candidate Tariq Iqbal mounted the final defense of his dissertation in the field of robotics on June 8. Iqbal followed his advisor, professor Laurel
Publications - Stanford University
AN INTEGRATED APPROACH TO ROBOTIC NAVIGATION UNDER UNCERTAINTY A DISSERTATION SUBMITTED TO THE DEPARTMENT OF ELECTRICAL in scope and quality as a dissertation for the degree of Doctor of Philosophy. Tze Lai, Primary Adviser want to thank Prof. Papanicolaou for being the Chair of my PhD oral defense. | http://transporqtu.cf/phd-thesis-robot-447010.html |
This magenta bike on a green background looks like a demonstration of the traditional colour wheel but it’s a real bike in a real street of Montmartre. Nov 2013.
This makes for an abstract picture: the cement dome of St Jean de Montmartre which I described in an earlier post is geometrical and beautiful in the sunshine.
Art Nouveau in Montmartre (2)
Another view of Saint Jean de Montmartre, the Art Nouveau church also nicknamed St Jean des briques because it is made of bricks and not stones as is usually the case there.
Getting to Yes!
Isn't that a positive shop window! Definitely so ... homies, Abbesses, Montmartre, February 2013
Art Nouveau in Montmartre
Art Nouveau is ubiquitous in Paris. Opposite the Abbesses metro station in Montmartre, there is this church, an austere yet beautifully crafted and picturesque church named Saint Jean de Montmartre. It was made of cement - not concrete - a special technique which prevailed until 1915 but was replaced with reinforced concrete later on, as with... Continue Reading →
Abbesses (Montmartre)
The abbesses metro station is undoubtedly and reputedly one of the most beautiful Art Nouveau buildings of the Paris underground train system. Its glass canopy is only one of two of those still available today. Yet, the station isn't the original one; its entrance was taken from another location and rebuilt here in the 1970s. | https://antimuseum.com/category/paris/montmartre/ |
Trash Andy Mulligan Essay Sample
- Pages:
- Word count: 459
- Category: prison
Get Full Essay
Get access to this section to get all help you need with your essay and educational issues.Get Access
Trash Andy Mulligan Essay Sample
Topic sentence quality convention
Examples/evidence context quotes
Explanation of examples
Concluding sentence
Olivia Weston is the temporary house mother at Behala’s Mission School and she has been characterised as a compassionate individual who wants to make a difference to the children’s lives. Olivia’s compassionate nature is revealed primarily through her thoughts and behaviour. Part way through the novel Olivia recounts her trip to Colva Prison with the boys. She begins this section explaining how she “fell in love” with the Behala children and the “eyes looking at me, and the smiles” (p.78). She goes on to share that visiting “the mountains of trash, and the children… is a thing to change your life” (p.78). Olivia’s thoughts immediately position the reader to understand the depth of her affection for the Behala children and her desire to care for them.
Her compassionate nature is further reinforced through her behaviour when she helps the boys visit the prison. In fact, not only does she act as their escort, she pays for their new clothes even though the “prices stunned [her]”, and she pays for the taxi fare even though she “gulped when [she] saw the meter” (p.83). Clearly, Olivia does whatever she can to help the boys, despite the fact that they achieve their goals at her expense.
Characterisation via thoughts and behaviour has positioned the reader to view Olivia as a compassionate individual, whose admirable qualities often result in her being manipulated by those she most cares for.
Topic sentence value agree/disagree
Examples/evidence context quotes
Explanation of examples
Examples/evidence transition
Explanation of examples
Concluding/linking sentence
One of the central values promoted in Trash, which strongly aligns with my own belief system, is ‘community’ and an appreciation for the support that communities offer. Valuing one’s community is presented through Raphael and his traumatic experience at the police station. When he returned to Behala the “whole neighbourhood came out” because when “one of their numbers is hurt, everyone feels the wound” (p.79). Raphael is grateful for the community’s care and compassion, and it helps him recover from the incident and continue to solve the Jose Angelico mystery. As a teacher, I constantly experience the benefits of belonging to a strong school community.
At the moment in WA, the state government is cutting funding from public education and teachers will be involved in industrial action. My school’s board has endorsed the teachers’ actions and requested that students and parents get behind us too. I am grateful for their support because we will not be successful if we do not have a united voice, and this speaks to why I value the idea of community. On this occasion, my values align with the values promoted in Trash, but this is not the case when it comes to trust. | https://blablawriting.com/trash-andy-mulligan-essay |
Teachers hate meetings, they just do. If you ask a teacher how a meeting went, they’d likely respond with, “it could’ve been an email.” We’ve all been there. I’m even guilty of giving that response. My theory as to why teachers hate meetings is because meetings take time. Teachers do not have a lot of extra time throughout the course of a school day, nearly every minute is accounted for. When meetings get scheduled, that takes time that was already allocated for something else the teacher had planned. The worst possible outcome for a meeting is for a teacher to feel like that allocation of their time was wasteful or inefficient.
As an instructional coach, you are already facing an uphill battle with teachers. Coaching requires meetings and meetings require teachers to reallocate some of their time, which can be overwhelming or burdensome on their daily schedule. I have had several teachers back out of instructional coaching because of the time commitment it requires. It is imperative that as instructional coaches we ensure that our coaching meetings are productive and efficient for our teachers. If your meetings are valued by your teachers, they will go out of their way to make them a priority. The first coaching meeting is the most essential. If you can win teachers over in the first meeting, they will be more open-minded and coachable for the subsequent meetings that will follow. I currently have several teachers on my caseload that tell me they look forward to our coaching conversations, and have asked to be coached for the duration of the school year. The best way to ensure that coaching conversations are viewed as a benefit to your teachers is to come to meetings prepared, with clear goals and objectives, and a plan.
As teachers, we were required to lesson plan for our classes. Over time, and with experience, you may have been able to survive a class period or two without a lesson plan, but lessons are smoother, more strategic and tend to dive deeper when thoughtfully planned ahead of time. Coaching conversations are no different. Just as teachers need to come to class with a lesson plan, instructional coaches must also plan – on paper – for their coaching conversations. One of the greatest mistakes a new coach can make is going into coaching meetings unprepared – I know, because I’ve done it.
When I first started coaching, I took the “I’ll go where the conversation takes me” approach, but what I’ve found is that it is an inefficient use of time. The first few minutes of the meeting were typically spent in awkward silence as the teacher prepared themselves for the conversation. Next came several minutes of pleasantries and ‘get to know you’ type questions. Then I would ask about what teachers are doing in their classes and what they’d like support with in the hopes that something would spark an idea. More often than not, the teacher would just go over their lesson plans, I’d say “that sounds great” and we would go our separate ways. Sometimes a teacher would ask me if I had a tool, strategy or activity for a particular class and I would have to say, “I don’t know, let me do some work and get back to you.” Had I prepared ahead of time, I could come to meetings with pointed questions, recommendations and resources, thus maximizing the time we had together.
So how do we create this plan?
Planning for a coaching conversation is similar in many ways to planning a lesson. Each lesson has clear goals, a route to meet those goals, and anticipated challenges that might arise. Coaching conversations should follow the same path. As you work with teachers, one of the first things you do is identify the challenge they want to work on. That challenge becomes the goal. Next, you identify strategies and tools to help them address that challenge. That becomes the route to meet the goal. Just as teachers need to think about challenges that might arise in their lessons, coaches need to consider challenges as well. Will a strategy work for different class sizes? Does the teacher have the resources they need? Does the teacher need to learn how to use a particular tool or program for a lesson?
These are all considerations the coach needs to make as they prepare for a coaching conversation. Teachers will often review material that might be helpful, practice the activity themselves, or even keep notes closeby as they teach to help guide the lesson when needed. Coaches should do the same. I rarely recommend a tool or a strategy that I haven’t already used myself, because as teachers implement something new they will have questions, and as an instructional coach, you need to be prepared to answer those or at least help troubleshoot. Similarly to planning a lesson, when you plan a coaching conversation you need to keep in mind that you may need to change course, modify plans, or even abandon what you had planned because some other pressing need presents itself throughout the conversation. So, let’s dive on in.
Step 1: Identify the objective for the conversation
Is this an initial coaching conversation where you need to explain how instructional coaching will work? Is this the goal-setting meeting in which you help the teacher identify their challenge area? Is this a follow-up meeting after a classroom visit or implementation of a new instructional tool or strategy?
The first step to plan a coaching conversation is to identify where the teacher needs to go in a particular meeting. To do this, you can read over the notes and reflections you made after the last coaching meeting. This helps you as the instructional coach remember what challenges your teacher is having, what strategies or tools you’ve already recommended and gives you a place to begin the conversation. For an instructional coach, keeping these types of notes is essential in planning for effective conversations. It does not matter if you take handwritten notes, use a collaborative Google Doc or even have something more robust like the DLP Coaching Dashboard, what is important is keeping track of each conversation you have with teachers. With several teachers on your caseload, you will not remember everything from meetings each week. By the same token, teachers will not remember everything discussed in your coaching meetings either. I recently met with a teacher and we discussed the use of digital rubrics, and then went through some ways she could organize the data that was collected. I praised the teacher for picking up on this skill so quickly, and she made a comment “yeah, but I have a class coming in soon so by the time I get back to working on this I’ll probably forget everything”. Teachers have so much on their plates already and it’s impossible for them to remember everything all the time. These notes allow both the teacher and coach to review where they’ve been and plan where they need to go. Think of it like planning for different classes, each one requires their own unique lesson plan based on where the class ended in the previous lesson.
Step 2: Identify tools and resources for scaffolding
The second thing to consider when planning for a coaching conversation is to think about where the teacher may need to go in order to move closer to their goals. Just as with students, teachers may need to have some scaffolding to get from where they started to where they want to end up. Each coaching meeting should move teachers a little bit further along than they were until the ultimate goal is reached. Consider this, when you were first learning to drive a car, did the instructor teach you to parallel parking, merge onto the highway and perform a 3 point turn all in the same lesson? No. Just as students often need material chunked and scaffolded, so do teachers. This also helps ensure that coaching meetings continue to be productive, because each one builds upon the meeting before it. Teachers are able to learn one skill or tool, become comfortable with it and then add more depth or complexity later on. This is very similar to working with students. As a teacher, you must constantly assess where your students are, where they need to be and what the next step is in getting them there.
Step 3: Script your meeting
Though planning for a coaching meeting is important, it is just as important that you don’t over plan. Remember, we do not want our coaching conversations to turn into coaching interrogations. When I first started teaching, I would script my entire class. I would write down exactly what I wanted to say at each step of the lesson and it would throw me for a loop when the student’s wouldn’t respond the way I had scripted it. Don’t worry, I learned how to adapt. Just as teaching, coaching conversations should be organic and ebb and flow where needed. Teachers, much like students, may not realize if we subtly guide the conversation in a particular direction with thoughtful questions, but as instructional coaches we also need to be prepared to go in a different direction entirely based on our teacher’s needs.
The number one barrier for many coaches is time. Time is an obstacle not just in ensuring meetings are productive, but it takes time to prepare for each meeting. As a teacher, lesson planning takes time. Teachers are often afforded a planning period, because schools recognize that planning effective lessons requires time and they want to ensure that teachers have the opportunity to do so. As an instructional coach, we don’t get a planning period. What we do get, however, is the ability to moderate our own schedules. Set aside a planning period for yourself each day or each week. Block off that time on your calendar to protect it, and use that time to review previous meeting notes and plan for upcoming meetings. Some conversations take a long time to plan for and some take less. The more coaching conversations you have, the more efficiently you’ll be able to plan.
There is an infinite amount of resources out there about teaching, education and instructional practices. As an instructional coach, we do our best to stay up-to-date with as much information as possible, but there is no guarantee we are going to know or remember all of it. We are more likely to guarantee that we won’t. As you plan your coaching conversations, you have an opportunity to gather some materials or resources you think may be relevant for each of your teachers. Even if the conversation digresses, you still have some tangible ideas you can leave with teachers or tools they can implement right away. Planning instructional coaching conversations allows you as an instructional coach to keep the big picture in mind and consider the needs of your teachers as they work toward their instructional goals.
Megan Purcell is a Digital Learning Specialist and Certified Dynamic Learning Project coach in Carrollton-Farmers Branch ISD located in Carrollton, TX. She enjoys working with teachers to help them elevate their teaching through the use of impactful technology tools and strategies. Megan holds a masters degree in Educational Technology, which she earned overseas at the National University of Ireland in Galway, in addition to being a certified Microsoft Innovative Educator and Apple Teacher. She is a former high school English teacher who loves learning, technology, and helping make life easier for her teachers. She believes that every student should have access to current technology in order to develop 21st century skills necessary for participating in a global society. | https://www.edtechteam.com/blog/2020/02/how-to-structure-an-effective-coaching-conversation/ |
Michael Blaustone is a certified Nutritional Therapy Practitioner and is certified through the Nutritional Therapy Association (NTA). Michael serves as not only as a consultant, but as a paraprofessional clinician utilizing hands on Functional Evaluation Testing to determine nutritional deficiencies in order to maximize support for his clients.
The NTA’s rigorous course curriculum requires proficiency in anatomy and physiology, basic chemistry concepts, and the science of food and its’ nutritional components. NTA Practitioners identify and address foundational imbalances in the body through intervention with nutrition in order to improve overall health and wellness.
After personally experiencing the healing power of nutrition and our body’s innate ability to utilize a nutrient dense diet to effectively manage and maintain optimal health, Michael shifted his educational studies from the standard of care currently provided by conventional medicine to that of Holistic Nutritional Therapy.
Michael’s passion is sharing his knowledge and skills in order to empower all of us to make smarter choices through lifestyle and a nutrient dense, whole food diet, in order to maximize our health, happiness and vitality.
Michael is also an active member of the Weston A. Price Foundation whose mission is dedicated to restoring nutrient-dense foods to the human diet through education, research and activism. | http://outlookonwellness.com/?page_id=49 |
Subsets and Splits
No community queries yet
The top public SQL queries from the community will appear here once available.