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Personal Trainer Guide
If you’re looking to learn more before you find a Trainer or, if you’ve decided that you might want to become one, this is the place for you to start.
What is a Personal Trainer?
A Trainer is a fitness professional who combines exercise science and experience to sherpa you towards your goals.
Have you been spending a lot of your time focus on other things, other people? If so, it might be time to make yourself a priority.
Is it necessary to hire a personal trainer? Nope, it's not.
That being said - If you weren't motivated, you wouldn't be here.
Let's channel that spark into a flame.
Personal Trainer Steven Mack Smiling in front of the Simple Solutions Fitness Studio in a blue logoed shirt
Written by Steven Mack, a Trainer located in Columbia, Missouri.
What You'll Learn Here:
Why Use A Personal Trainer?
• Why should you hire a trainer? Motivation? Exercises?
• Does training work?
• What are some of the disadvantages?
Talking cartoon icon with blue bubble
What Do Personal Trainers Do?
• What services do trainers offer?
• Can your trainer give you a massage? Help you stretch? Provide meal plans?
• Where do personal trainers work?
an envelope with a green dollar bill sticking out of it
How Much Does a Personal Trainer Cost?
• What should you expect to pay for a trainer per month? 12 weeks?
• Why don't gyms list prices?
• Why are personal trainers so expensive? What does a trainer make?
How to Find a Personal Trainer
• How do you know if your trainer is good? Does your trainer have to be fit?
• What should you ask your trainer?
• How to find a good personal trainer, when you're ready,
How Do You Become a Personal Trainer?
• What do gyms look for in a personal trainer?
• How hard is it to become a personal trainer? How much can you expect to make?
• Why do trainers leave the field?
How Long Should You Have a Personal Trainer?
• How long does it take to see results?
• How fit can you get in 3 months?
• Also related to how much personal training costs.
Ready to talk to a human?
Reach out by filling out the contact form. You'll get a call the moment it's seen.
Related Questions People Also Ask About Training:
What Does a Personal Trainer Do for You?
A pt can help guide you towards your fitness goals of feeling strong and confident.
Each trainer will have their procedures but, the process might look something like the following.
After you sit down for an initial consultation, you’ll both be clear on:
• What your goals are.
• How you’re going to workout each session.
• How you’re going to measure your progress (or not if you don’t like scales).
• When you’ll be showing up to workout.
• What you’re going to change when you eventually reach a plateau.
After you get clear on your goals, budget, and some logistics, you’ll meet with your trainer to carry out your fitness plan. Most trainers can give you general nutrition advice or will refer out when necessary.
Do You Really Need a Personal Trainer?
You don’t need a pt if you’re comfortable teaching yourself a skill.
You can learn to workout using:
• YouTube videos
• Step-by-step exercise guides and books
• Strength training programs (online/offline)
Over time, your results will compound. If you're making consistent progress, it's all a matter of staying in the game.
You might want to hire a trainer if you feel that you would benefit from socialization or having someone Sherpa you along your journey.
Your rate of progress will slow down with strength training. The key will be in knowing what you need to try next to keep making progress.
Explainer articles like these four can help get you started:
What Are the Qualities of a Good Trainer?
A fitness trainer needs to have the ability to ask questions, listen, and communicate.
You want to feel like a human is listening to you and considering your needs.
Some basic functions of a training business include:
• Being able to demonstrate and coach exercises.
• Walking clients through nutritional roadblocks.
• Keeping a space tidy and clean for training.
• Remembering personal details about your life.
Once these basic expectations are met, you can look to make things more personal to each client.
This includes modifying exercise routines, using names, asking for feedback, and celebrating milestones.
How Long is Personal Training?
Most pts work with clients for 30, 45, or 60 minutes. If you’re an endurance athlete or play a sport, a session may be as long as 90 or 120 minutes.
The length of a workout depends on the way that your program and schedule are laid out.
Let’s say that you’ve figured out that you want to build strong legs. There are many ways to follow a routine aimed at that goal.
For example:
• Will you workout twice a week, three times, or six?
• Are you going to perform the same exercises every day or change movements every day?
• Will you be working out for 30 minutes each time or an hour?
• Is your schedule pretty consistent or do you need a more flexible style of training?
Each of those answers can affect how long your training session lasts.
You could train four times a week for 30 minutes or two times a week for 60 and still do the same total amount of exercise volume.
Can a Personal Trainer Bill Insurance?
Yes, a trainer can bill insurance if it is prescribed by your primary care physician.
I have worked with a few clients who pay for training using their HSA after receiving the all-clear.
Talk to your healthcare provider to find out if this is something you have the option of pursuing.
Do Personal Trainers Work Out with Their Clients?
Fitness trainers do not work out with their clients.
Trainers supervise clients as they perform exercises.
A trainer is looking to:
• Explain and demonstrate an exercise before it is performed.
• Watch your form during a set, providing cues as needed.
• Provide informational feedback after a set or session.
• Count the number of repetitions you perform (sometimes).
Group fitness instructors may workout with their clients. Some trainers work as fitness instructors. Fitness instructors are not trainers.
Are You Supposed to Tip Your Personal Trainer?
You don’t need to tip your trainer. It would be incredibly helpful if you gave your trainer some feedback.
Feedback can be worth its weight in gold if you’re willing to do some listening.
Often, it can be really hard for businesses to get feedback from customers. You’re the expert on your body. Your trainer can do a better job of helping you when you take the time to talk about yourself and how things are going.
How Much is Too Much for a Personal Trainer?
If you cannot work with a trainer for more than a few months, you might consider another fitness option. To benefit from exercise, you will need to work out with some consistency and frequency.
Aim to lift weights at least twice a week for four to twelve or so weeks.
In a month, you might notice some changes - nothing that will break the internet. Over time, your results will compound, if you’re on the right track.
You might find a trainer that is willing to be flexible with you on scheduling. Some trainers write programs and check in with clients as frequently as they see fit.
Red headed girl wearing glasses staring intently at a macbook laptop while laying in bed
Do You Need a Degree to Be a Personal Trainer?
You do not need a degree to become a fitness trainer.
The first step to becoming a trainer is getting some experience. One of the most interesting things about this industry is the sheer number of paths that people have taken to get to it.
Training can be a first, second, or part-time career for some of the most influential trainers in the market.
I had the pleasure of contributing to two articles you may want to read about in-person and online training:
In the interviews I conducted for the online training article, almost every trainer mentioned a different origin story:
Learn more about getting started as a personal trainer in this guide.
Do You Need a Certification to be a Personal Trainer?
Can you be a personal trainer without a certification? Yes.
Personal training does not require a license or certification. While you don’t need a certification to become a trainer, it’s viewed as a liability to go without.
Depending on the cert, your trainer may or may not have an educational standard.
Expectations for trainers also vary depending on who they work for (or don’t).
A college degree is beyond the educational requirements of most training certifications.
The National Strength and Conditioning Association (NSCA) requires that trainers:
1. Be at least 18 years old
2. Have a high school diploma or equivalent.
3. Have a current CPR/AED certification.
(I suspect COVID-19 will have some effect on CPR/AED certification requirement)
Which Personal Trainer Certification is Best?
The NSCA Certified Strength and Conditioning Specialist Certification is considered the best.
The CSCS is hard to get and requires a bachelor's degree - just to take the test.
In actuality, there isn’t a way to objectively compare personal training certifications. I’ve held a couple, right now a CSCS. I firmly believe in the NSCAs mission.
The NSCA offers a personal training certification, CPT. It seems to be easier to test for than the CSCS. In 2018, 58% and 78% of candidates passed the CSCS and CPT exams respectively.
While I have no experience with it, the ACSM-CPT carries a high level of regard within the industry as well.
Want to Start a Training Program?
If you didn’t learn to lift in high school or it’s been a while, a weight room looks like a maze of weights, switches, knobs and bros.
Do more than just go pick up a dumbell; let’s create a plan to build strength, get sustainable results (and look like you know what you’re doing).
We’ll start with the basics of strength training and you’ll walk out with (a starter amount of) confidence.
Phone: 1-573-443-1495
• I want to see changes in my:
• I want to start:
• I am interested in training:
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电解抛光后铝箔表面粗糙度对AAO模板孔间距的影响
Interpore distance dependence of anodic aluminium oxide films on high-purity aluminium sheet surface roughness controlled by electropolishing
姜永军，许刚茜，王亚兰
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作者单位：内蒙古科技大学机械学院
中文关键字： AAO模板；粗糙度；孔间距；电场
英文关键字：AAO membranes; surface roughness; interpore distances; electric field
中文摘要：利用不同粗糙度的高纯铝箔在0.3M草酸中制备阳极氧化铝模板（AAO），研究发现随粗糙度由5到50nm增大，孔间距从80nm减小至70nm。采用有限元数值分析（FEM）方法建立二维模型，深入研究粗糙度对氧化电场的影响。结果表明铝表面的显微不平造成了反向电场。一定范围内，粗糙度增大，反向电场增强，膜内的平均电场被削弱，进而制备的模板孔间距减小。
英文摘要：Self-organized nano-porous anodic aluminium oxide (AAO) membranes were obtained on the high purity aluminum foil with different surface roughness after electropolishing in 0.3M oxalic acid solution.The average interpore distances of the as-prepared AAO membranes reduces from 80nm to 70nm as the surface roughness increases from 5nm to 50nm. Three two-dimensional models of alumina with different roughness was constructed to calculate the electric field distributions based on Finite Element Method (FEM). The result show that the micro rough on the foil samples generated a reverse electric field . The reverse electric field could be increased and the average electric field within the membrane could be weakened by increasing the roughness, resulting in the nano interpore distance of template is reduced.
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"Fossies" - the Fresh Open Source Software Archive
Source code changes of the file "poetry/console/commands/env/list.py" between
poetry-1.1.15.tar.gz and poetry-1.2.0.tar.gz
About: Poetry is a tool for dependency management and packaging in Python.
list.py (poetry-1.1.15):list.py (poetry-1.2.0)
from cleo import option from __future__ import annotations
from ..command import Command from cleo.helpers import option
class EnvListCommand(Command): from poetry.console.commands.command import Command
name = "list" class EnvListCommand(Command):
name = "env list"
description = "Lists all virtualenvs associated with the current project." description = "Lists all virtualenvs associated with the current project."
options = [option("full-path", None, "Output the full paths of the virtualen vs.")] options = [option("full-path", None, "Output the full paths of the virtualen vs.")]
def handle(self): def handle(self) -> int:
from poetry.utils.env import EnvManager from poetry.utils.env import EnvManager
manager = EnvManager(self.poetry) manager = EnvManager(self.poetry)
current_env = manager.get() current_env = manager.get()
for venv in manager.list(): for venv in manager.list():
name = venv.path.name name = venv.path.name
if self.option("full-path"): if self.option("full-path"):
name = str(venv.path) name = str(venv.path)
if venv == current_env: if venv == current_env:
self.line("<info>{} (Activated)</info>".format(name)) self.line(f"<info>{name} (Activated)</info>")
continue continue
self.line(name) self.line(name)
return 0
End of changes. 7 change blocks.
6 lines changed or deleted 7 lines changed or added
Home | About | Features | All | Newest | Dox | Diffs | RSS Feeds | Screenshots | Comments | Imprint | Privacy | HTTP(S)
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// !! IMPORTANT README: // You may add additional external JS and CSS as needed to complete the project, however the current external resource MUST remain in place for the tests to work. BABEL must also be left in place. /*********** INSTRUCTIONS: - Select the project you would like to complete from the dropdown menu. - Click the "RUN TESTS" button to run the tests against the blank pen. - Click the "TESTS" button to see the individual test cases. (should all be failing at first) - Start coding! As you fulfill each test case, you will see them go from red to green. - As you start to build out your project, when tests are failing, you should get helpful errors along the way! ************/ // PLEASE NOTE: Adding global style rules using the * selector, or by adding rules to body {..} or html {..}, or to all elements within body or html, i.e. h1 {..}, has the potential to pollute the test suite's CSS. Try adding: * { color: red }, for a quick example! // Once you have read the above messages, you can delete all comments. "use strict";
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Nuclear Radiation and Health Effects
Manostaxx
The text that follows is owned by the site above referred.
Here is only a small part of the article, for more please follow the link
SOURCE: http://www.world-nuclear.org/information-library/safety-and-security/radiation-and-health/nuclear-radiation-and-health-effects.aspx
• Natural sources account for most of the radiation we all receive each year.
• The nuclear fuel cycle does not give rise to significant radiation exposure for members of the public, and even in two major nuclear accidents – Three Mile Island and Fukushima – exposure to radiation has caused no harm to the public.
• Radiation protection standards assume that any dose of radiation, no matter how small, involves a possible risk to human health. This deliberately conservative assumption is increasingly being questioned.
• Fear of radiation causes much harm. Expressed particularly in government edicts following the Fukushima accident (and also Chernobyl), it has caused much suffering and many deaths.
Radiation is energy in the process of being transmitted. It may take such forms as light, or tiny particles much too small to see. Visible light, the ultra-violet light we receive from the sun, and transmission signals for TV and radio communications are all forms of radiation that are common in our daily lives. These are all generally referred to as ‘non-ionizing’ radiation, though at least some ultra-violet radiation is considered to be ionizing.
Radiation particularly associated with nuclear medicine and the use of nuclear energy, along with X-rays, is ‘ionizing’ radiation, which means that the radiation has sufficient energy to interact with matter, especially the human body, and produce ions, i.e. it can eject an electron from an atom.
X-rays from a high-voltage discharge were discovered in 1895, and radioactivity from the decay of particular isotopes was discovered in 1896. Many scientists then undertook study of these, and especially their medical applications. This led to the identification of different kinds of radiation from the decay of atomic nuclei, and understanding of the nature of the atom. Neutrons were identified in 1932, and in 1939 atomic fission was discovered by irradiating uranium with neutrons. This led on to harnessing the energy released by fission.
A 2016 United Nations Environment Programme (UNEP) publication notes: “Today, we know more about the sources and effects of exposure to [ionizing] radiation than to almost any other hazardous agent, and the scientific community is constantly updating and analysing its knowledge… The sources of radiation causing the greatest exposure of the general public are not necessarily those that attract the most attention. In fact, the greatest exposure is caused by natural sources ever present in the environment, and the major contributor to exposure from artificial sources is the use of radiation in medicine worldwide.”
Types of radiation
Nuclear radiation arises from hundreds of different kinds of unstable atoms. While many exist in nature, the majority are created in nuclear reactionsa. Ionizing radiation which can damage living tissue is emitted as the unstable atoms (radionuclides) change (‘decay’) spontaneously to become different kinds of atoms.
The principal kinds of ionizing radiation are:
Alpha particles
These are helium nuclei consisting of two protons and two neutrons and are emitted from naturally-occurring heavy elements such as uranium and radium, as well as from some man-made transuranic elements. They are intensely ionizing but cannot penetrate the skin, so are dangerous only if emitted inside the body.
Beta particles
These are fast-moving electrons emitted by many radioactive elements. They are more penetrating than alpha particles, but easily shielded – the most energetic of them can be stopped by a few millimetres of wood or aluminium. They can penetrate a little way into human flesh but are generally less dangerous to people than gamma radiation. Exposure produces an effect like sunburn, but which is slower to heal. The weakest of them, such as from tritium, are stopped by skin or cellophane. Beta-radioactive substances are also safe if kept in appropriate sealed containers.
Gamma rays
These are high-energy electromagnetic waves much the same as X-rays. They are emitted in many radioactive decays and may be very penetrating, so require more substantial shielding. Gamma rays are the main hazard to people dealing with sealed radioactive materials used, for example, in industrial gauges and radiotherapy machines. Radiation dose badges are worn by workers in exposed situations to detect them and hence monitor exposure. All of us receive about 0.5-1 mSv per year of gamma radiation from rocks, and in some places, much more. Gamma activity in a substance (e.g. rock) can be measured with a scintillometer or Geiger counter.
X-rays are also electromagnetic waves and ionizing, virtually identical to gamma rays, but not nuclear in origin. They are produced in a vacuum tube where an electron beam from a cathode is fired at target material comprising an anode, so are produced on demand rather than by inexorable physical processes. (However the effect of this radiation does not depend on its origin but on its energy.)
Cosmic radiation consists of very energetic particles, mostly high-energy protons, which bombard the Earth from outer space. They comprise about one-tenth of natural background exposure at sea level, and more at high altitudes.
Neutrons are uncharged particles mostly released by nuclear fission (the splitting of atoms in a nuclear reactor), and hence are seldom encountered outside the core of a nuclear reactor.* Thus they are not normally a problem outside nuclear plants. Fast neutrons can be very destructive to human tissue. Neutrons are the only type of radiation which can make other, non-radioactive materials, become radioactive.
* Large nuclei can fission spontaneously, since the so-called strong nuclear force holding each nucleus together is not overwhelmingly stronger than the repulsive force of charged protons.
Units of radiation and radioactivity
In order to quantify how much radiation we are exposed to in our daily lives and to assess potential health impacts as a result, it is necessary to establish a unit of measurement. The basic unit of radiation dose absorbed in tissue is the gray (Gy), where one gray represents the deposition of one joule of energy per kilogram of tissue.
However, since neutrons and alpha particles cause more damage per gray than gamma or beta radiation, another unit, the sievert (Sv) is used in setting radiological protection standards. This weighted unit of measurement takes into account biological effects of different types of radiation and indicates the equivelent dose. One gray of beta or gamma radiation has one sievert of biological effect, one gray of alpha particles has 20 Sv effect and one gray of neutrons is equivalent to around 10 Sv (depending on their energy). Since the sievert is a relatively large value, dose to humans is normally measured in millisieverts (mSv), one-thousandth of a sievert.
Note that Sv and Gy measurements are accumulated over time, whereas damage (or effect) depends on the actual dose rate, e.g. mSv per day or year, Gy per day in radiotherapy.
The becquerel (Bq) is a unit or measure of actual radioactivity in material (as distinct from the radiation it emits, or the human dose from that), with reference to the number of nuclear disintegrations per second (1 Bq = 1 disintegration/sec). Quantities of radioactive material are commonly estimated by measuring the amount of intrinsic radioactivity in becquerels – one Bq of radioactive material is that amount which has an average of one disintegration per second, i.e. an activity of 1 Bq. This may be spread through a very large mass.
Radioactivity of some natural and other materials
1 adult human (65 Bq/kg) 4500 Bq
1 kg of coffee 1000 Bq
1 kg of brazil nuts 400 Bq
1 banana 15 Bq
The air in a 100 sq metre Australian home (radon) 3000 Bq
The air in many 100 sq metre European homes (radon) Up to 30 000 Bq
1 household smoke detector (with americium) 30 000 Bq
Radioisotope for medical diagnosis 70 million Bq
Radioisotope source for medical therapy 100 000 000 million Bq (100 TBq)
1 kg 50-year old vitrified high-level nuclear waste 10 000 000 million Bq (10 TBq)
1 luminous Exit sign (1970s) 1 000 000 million Bq (1 TBq)
1 kg uranium 25 million Bq
1 kg uranium ore (Canadian, 15%) 25 million Bq
1 kg uranium ore (Australian, 0.3%) 500 000 Bq
1 kg low level radioactive waste 1 million Bq
1 kg of coal ash 2000 Bq
1 kg of granite 1000 Bq
1 kg of superphosphate fertilizer 5000 Bq
N.B. Though the intrinsic radioactivity is the same, the radiation dose received by someone handling a kilogram of high-grade uranium ore will be much greater than for the same exposure to a kilogram of separated uranium, since the ore contains a number of short-lived decay products (see section on Radioactive Decay), while the uranium has a very long half-life.
Older units of radiation measurement continue in use in some literature:
1 gray = 100 rads
1 sievert = 100 rem
1 becquerel = 27 picocuries or 2.7 x 10-11 curies
One curie was originally the activity of one gram of radium-226, and represents 3.7 x 1010 disintegrations per second (Bq).
The Working Level Month (WLM) has been used as a measure of dose for exposure to radon and in particular, radon decay productsb.
Since there is radioactivity in many foodstuffs, there has been a whimsical suggestion that the Banana Equivalent Dose from eating one banana be adopted for popular reference. This is about 0.0001 mSv.
Routine sources of radiation
Radiation can arise from human activities or from natural sources. Most radiation exposure is from natural sources. These include: radioactivity in rocks and soil of the Earth’s crust; radon, a radioactive gas given out by many volcanic rocks and uranium ore; and cosmic radiation. The human environment has always been radioactive and accounts for up to 85% of the annual human radiation dose.
Helpful depictions of routine sources of radiation can be found on the information is beautiful and xkcd websites.
Radiation arising from human activities typically accounts for up to 20% of the public’s exposure every year as global average. In the USA by 2006 it averaged about half of the total. This radiation is no different from natural radiation except that it can be controlled. X-rays and other medical procedures account for most exposure from this quarter. Less than 1% of exposure is due to the fallout from past testing of nuclear weapons or the generation of electricity in nuclear, as well as coal and geothermal, power plants.
Backscatter X-ray scanners being introduced for airport security will gives exposure of up to 5 microsieverts (μSv), compared with 5 μSv on a short flight and 30 μSv on a long intercontinental flight across the equator, or more at higher latitudes – by a factor of 2 or 3. Aircrew can receive up to about 5 mSv/yr from their hours in the air, while frequent flyers can score a similar incrementc. On average, nuclear power workers receive a lower annual radiation dose than flight crew, and frequent flyers in 250 hours would receive 1 mSv.
The maximum annual dose allowed for radiation workers is 20 mSv/yr, though in practice, doses are usually kept well below this level. In comparison, the average dose received by the public from nuclear power is 0.0002 mSv/yr, which is of the order of 10,000 times smaller than the total yearly dose received by the public from background radiation.
Sources of Radiation pie graph
Natural background radiation, radon
Naturally occurring background radiation is the main source of exposure for most people, and provides some perspective on radiation exposure from nuclear energy. Much of it comes from primordial radionuclides in the Earth’s crust, and materials from it. Potasssium-40, uranium-238 and thorium-232 with their decay products are the main source.
The average dose received by all of us from background radiation is around 2.4 mSv/yr, which can vary depending on the geology and altitude where people live – ranging between 1 and 10 mSv/yr, but can be more than 50 mSv/yr. The highest known level of background radiation affecting a substantial population is in Kerala and Madras states in India where some 140,000 people receive doses which average over 15 millisievert per year from gamma radiation, in addition to a similar dose from radon. Comparable levels occur in Brazil and Sudan, with average exposures up to about 40 mSv/yr to many people. (The highest level of natural background radiation recorded is on a Brazilian beach: 800 mSv/yr, but people don’t live there.)
Several places are known in Iran, India and Europe where natural background radiation gives an annual dose of more than 100 mSv to people and up to 260 mSv (at Ramsar in Iran, where some 200,000 people are exposed to more than 10 mSv/yr). Lifetime doses from natural radiation range up to several thousand millisievert. However, there is no evidence of increased cancers or other health problems arising from these high natural levels. The millions of nuclear workers that have been monitored closely for 50 years have no higher cancer mortality than the general population but have had up to ten times the average dose. People living in Colorado and Wyoming have twice the annual dose as those in Los Angeles, but have lower cancer rates. Misasa hot springs in western Honshu, a Japan Heritage site, attracts people due to having high levels of radium (up to 550 Bq/L), with health effects long claimed, and in a 1992 study the local residents’ cancer death rate was half the Japan average.* (Japan J.Cancer Res. 83,1-5, Jan 1992) A study on 3000 residents living in an area with 60 Bq/m3 radon (about ten times normal average) showed no health difference. Hot springs in China have levels reaching 3270 Bq/L radon-222 (Liaoning sanatorium), 2720 Bq/L (Tanghe hot spring) and 230 Bq/L (Puxzhe hot spring), though associated exposure from airborne radon are low**.
* The waters are promoted as boosting the body’s immunity and natural healing power, while helping to relieve bronchitis and diabetes symptoms, as well as beautifying the skin. Drinking the water is also said to have antioxidant effects. (These claims are not known to be endorsed by any public health authority.)
**
Chinese figures from Liu & Pan in NORM VII.
Radon is a naturally occurring radioactive gas resulting from the decay of uranium-238, which concentrates in enclosed spaces such as buildings and underground mines, particularly in early uranium mines where it sometimes became a significant hazard before the problem was understood and controlled by increased ventilation. Radon has decay products that are short-lived alpha emitters and deposit on surfaces in the respiratory tract during the passage of breathing air. At high radon levels, this can cause an increased risk of lung cancer, particularly for smokers. (Smoking itself has a very much greater lung cancer effect than radon.) People everywhere are typically exposed to around 0.2 mSv/yr, and often up to 3 mSv/yr, due to radon (mainly from inhalation in their homes) without apparent ill-effectd. Where deemed necessary, radon levels in buildings and mines can be controlled by ventilation, and measures can be taken in new constructions to prevent radon from entering buildings.
However, radon levels of up to 3700 Bq/m3 in some dwellings at Ramsar in Iran have no evident ill-effect. Here, a study (Mortazavi et al, 2005) showed that the highest lung cancer mortality rate was where radon levels were normal, and the lowest rate was where radon concentrations in dwellings were highest. The ICRP recommends keeping workplace radon levels below 300 Bq/m3, equivalent to about 10 mSv/yr. Above this, workers should be considered as occupationally exposed, and subject to the same monitoring as nuclear industry workers. The normal indoor radon concentration ranges from 10 to 100 Bq/m3, but may naturally reach 10,000 Bq/m3, according to UNEP.
Public exposure to natural radiatione
Source of exposure Annual effective dose (mSv)
Average Typical range
Cosmic radiation Directly ionizing and photon component 0.28
Neutron component 0.10
Cosmogenic radionuclides 0.01
Total cosmic and cosmogenic 0.39 0.3–1.0e
External terrestrial radiation Outdoors 0.07
Indoors 0.41
Total external terrestrial radiation 0.48 0.3-1.0e
Inhalation Uranium and thorium series 0.006
Radon (Rn-222) 1.15
Thoron (Rn-220) 0.1
Total inhalation exposure 1.26 0.2-10e
Ingestion K-40 0.17
Uranium and thorium series 0.12
Total ingestion exposure 0.29 0.2-1.0e
Total 2.4 1.0-13
Crews of nuclear submarines have possibly the lowest radiation exposure of anyone, despite living within a few metres of a nuclear reactor, since they are exposed to less natural background radiation than the rest of us, and the reactor compartment is well shielded.1 US Naval Reactors’ average annual occupational exposure was 0.06 mSv per person in 2013, and no personnel have exceeded 20 mSv in any year in the 34 years to then. The average occupational exposure of each person monitored at Naval Reactors’ facilities since 1958 is 1.03 mSv per year.
Effects of ionizing radiation
Some of the ultraviolet (UV) radiation from the sun is considered ionizing radiation, and provides a starting point in considering its effects. Sunlight UV is important in producing vitamin D in humans, but too much exposure produces sunburn and, potentially, skin cancer. Skin tissue is damaged, and that damage to DNA may not be repaired properly, so that over time, cancer develops and may be fatal. Adaptation from repeated low exposure can decrease vulnerability. But exposure to sunlight is quite properly sought after in moderation, and not widely feared.
Our knowledge of the effects of shorter-wavelength ionizing radiation from atomic nuclei derives primarily from groups of people who have received high doses. The main difference from UV radiation is that beta, gamma and X-rays can penetrate the skin. The risk associated with large doses of this ionizing radiation is relatively well established. However, the effects, and any risks associated with doses under about 200 mSv, are less obvious because of the large underlying incidence of cancer caused by other factors. Benefits of lower doses have long been recognised, though radiation protection standards assume that any dose of radiation, no matter how small, involves a possible risk to human health. However, available scientific evidence does not indicate any cancer risk or immediate effects at doses below 100 mSv per year. At low levels of exposure, the body’s natural mechanisms usually repair radiation damage to DNA in cells soon after it occurs (see following section on low-level radiation). However, high-level irradiation overwhelms those repair mechanisms and is harmful. Dose rate is as important as overall dose.
The UN Scientific Commission on the Effects of Atomic Radiation (UNSCEAR) currently uses the term low dose to mean absorbed levels below 100 mGy but greater than 10 mGy, and the term very low dose for any levels below 10 mGy. High absorbed dose is defined as more than about 1000 mGy. For beta and gamma radiation, these figures can be taken as mSv equivalent dose.
Some comparative whole-body radiation doses and their effects
2.4 mSv/yr Typical background radiation experienced by everyone (average 1.5 mSv in Australia, 3 mSv in North America).
1.5 to 2.5 mSv/yr Average dose to Australian uranium miners and US nuclear industry workers, above background and medical.
Up to 5 mSv/yr Typical incremental dose for aircrew in middle latitudes.
9 mSv/yr Exposure by airline crew flying the New York – Tokyo polar route.
10 mSv/yr Maximum actual dose to Australian uranium miners.
10 mSv Effective dose from abdomen & pelvis CT scan.
20 mSv/yr Current limit (averaged) for nuclear industry employees and uranium miners. (In Japan: 5 mSv per three months for women)
50 mSv/yr Former routine limit for nuclear industry employees, now maximum allowable for a single year (average to be 20 mSv/yr max). It is also the dose rate which arises from natural background levels in several places in Iran, India and Europe.
50 mSv Allowable short-term dose for emergency workers (IAEA).
100 mSv Lowest annual level at which increase in cancer risk is evident (UNSCEAR). Above this, the probability of cancer occurrence (rather than the severity) is assumed to increase with dose. No harm has been demonstrated below this dose.
Allowable short-term dose for emergency workers taking vital remedial actions (IAEA).
Dose from four months on international space station orbiting 350 km up.
130 mSv/yr Long-term safe level for public after radiological incident, measured 1 m above contaminated ground, calculated from published hourly rate x 0.6. Risk too low to justify any action below this (IAEA).
170 mSv/wk 7-day provisionally safe level for public after radiological incident, measured 1 m above contaminated ground (IAEA).
250 mSv Allowable short-term dose for workers controlling the 2011 Fukushima accident, set as emergency limit elsewhere.
250 mSv/yr Natural background level at Ramsar in Iran, with no identified health effects (Some exposures reach 700 mSv/yr). Maximum allowable annual dose in emergency situations in Japan (NRA).
350 mSv/lifetime Criterion for relocating people after Chernobyl accident.
500 mSv Allowable short-term dose for emergency workers taking life-saving actions (IAEA).
680 mSv/yr Tolerance dose level allowable to 1955 (assuming gamma, X-ray and beta radiation).
700 mSv/yr Suggested threshold for maintaining evacuation after nuclear accident.
(IAEA has 880 mSv/yr over one month as provisionally safe.
800 mSv/yr Highest level of natural background radiation recorded, on a Brazilian beach.
1,000 mSv short-term Assumed to be likely to cause a fatal cancer many years later in about 5 of every 100 persons exposed to it (i.e. if the normal incidence of fatal cancer were 25%, this dose would increase it to 30%).
Highest reference level recommended by ICRP for rescue workers in emergency situation.
1,000 mSv short-term Threshold for causing (temporary) radiation sickness (Acute Radiation Syndrome) such as nausea and decreased white blood cell count, but not death. Above this, severity of illness increases with dose.
5,000 mSv short-term Would kill about half those receiving it as whole body dose within a month. (However, this is only twice a typical daily therapeutic dose applied to a very small area of the body over 4 to 6 weeks or so to kill malignant cells in cancer treatment.)
10,000 mSv short-term Fatal within a few weeks.
The main expert body on radiation effects is the UN Scientific Commission on the Effects of Atomic Radiation (UNSCEAR), set up in 1955 and reporting to the UN General Assembly. It involves scientists from over 20 countries and publishes its findings in major reports. The UNSCEAR 2006 report dealt broadly with the Effects of Ionizing Radiation. Another valuable report, titled Low-level Radiation and its Implications for Fukushima Recovery, was published in June 2012 by the American Nuclear Society.
In 2012 UNSCEAR reported to the UN General Assembly on radiation effects. It had been asked in 2007 “to clarify further the assessment of potential harm owing to chronic low-level exposures among large populations and also the attributability of health effects” to radiation exposure. It said that while some effects from high acute doses were clear, others including hereditary effects in human populations were not, and could not be attributed to exposure, and that this was especially true at low levels. “In general, increases in the incidence of health effects in populations cannot be attributed reliably to chronic exposure to radiation at levels that are typical of the global average background levels of radiation.” Furthermore, multiplying very low doses by large numbers of individuals does not give a meaningful result regarding health effects. UNSCEAR also addressed uncertainties in risk estimation relating to cancer, particularly the extrapolations from high-dose to low-dose exposures and from acute to chronic and fractionated exposures. Earlier (1958) UNSCEAR data for leukaemia incidence among Hiroshima survivors suggested a threshold of about 400 mSv for harmful effects.
Average Annual Radiation Doses
Epidemiological studies continue on the survivors of the atomic bombing of Hiroshima and Nagasaki, involving some 76,000 people exposed at levels ranging up to more than 5,000 mSv. These have shown that radiation is the likely cause of several hundred deaths from cancer, in addition to the normal incidence found in any populationg. From this data the International Commission on Radiological Protection (ICRP) and others estimate the fatal cancer risk as 5% per sievert exposure for a population of all ages – so one person in 100 exposed to 200 mSv could be expected to develop a fatal cancer some years later. In Western countries, about a quarter of people die from cancers, with smoking, dietary factors, genetic factors and strong sunlight being among the main causes. About 40% of people are expected to develop cancer during their lifetime even in the absence of radiation exposure beyond normal background levels. Radiation is a weak carcinogen, but undue exposure can certainly increase health risks.
In 1990, the US National Cancer Institute (NCI) found no evidence of any increase in cancer mortality among people living near to 62 major nuclear facilities. The NCI study was the broadest of its kind ever conducted and supported similar studies conducted elsewhere in the USA as well as in Canada and Europe.h
About 60 years ago it was discovered that ionizing radiation could induce genetic mutations in fruit flies. Intensive study since then has shown that radiation can similarly induce mutations in plants and test animals. However there is no evidence of inherited genetic damage to humans from radiation, even as a result of the large doses received by atomic bomb survivors in Japan.
In a plant or animal cell the material (DNA) which carries genetic information necessary to cell development, maintenance and division is the critical target for radiation. Much of the damage to DNA is repairable, but in a small proportion of cells the DNA is permanently altered. This may result in death of the cell or development of a cancer, or in the case of cells forming gonad tissue, alterations which continue as genetic changes in subsequent generations. Most such mutational changes are deleterious; very few can be expected to result in improvements.
The relatively low levels of radiation allowed for members of the public and for workers in the nuclear industry are such that any increase in genetic effects due to nuclear power will be imperceptible and almost certainly non-existent. Radiation exposure levels are set so as to prevent tissue damage and minimize the risk of cancer. Experimental evidence indicates that cancers are more likely than inherited genetic damage.
Some 75,000 children born of parents who survived high radiation doses at Hiroshima and Nagasaki in 1945 have been the subject of intensive examination. This study confirms that no increase in genetic abnormalities in human populations is likely as a result of even quite high doses of radiation. Similarly, no genetic effects are evident as a result of the Chernobyl accident.
Life on Earth commenced and developed when the environment was certainly subject to several times as much radioactivity as it is now, so radiation is not a new phenomenon. If there is no dramatic increase in people’s general radiation exposure, there is no evidence that health or genetic effects from radiation could ever become significant.
Low-level radiation effects
Linear Hypothesis
A lot of research has been undertaken on the effects of low-level radiation. The findings have failed to support the so-called linear no-threshold (LNT) hypothesis. This theory assumes that the demonstrated relationships between radiation dose and adverse effects at high levels of exposure also applies to low levels and provides the (deliberately conservative) basis of occupational health and other radiation protection standards.
Increasing evidence suggests that there may be a threshold between 100 and 700 mSv below which no harmful effects of radiation occur without effective cellular repair. However, this is not yet accepted by national or international radiation protection bodies as sufficiently well-proven to be taken into official standards. Nevertheless, at low levels of exposure, the body’s natural mechanisms do repair radiation and other damage to cells soon after it occurs, and some adaptive response is stimulated which protects cells and tissues, as with exposure to other external agents at low levels. The ICRP recommends that the LNT model should be assumed for the purpose of optimising radiation protection practices, but that it should not be used for estimating the health effects of exposures to small radiation doses received by large numbers of people over long periods of time.
A November 2009 technical report from the Electric Power Research Institute in USA drew upon more than 200 peer-reviewed publications on effects of low-level radiation and concluded that the effects of low dose-rate radiation are different and that “the risks due to [those effects] may be over-estimated” by the linear hypothesis3. “From an epidemiological perspective, individual radiation doses of less than 100 mSv in a single exposure are too small to allow detection of any statistically significant excess cancers in the presence of naturally occurring cancers. The doses received by nuclear power plant workers fall into this category because exposure is accumulated over many years, with an average annual dose about 100 times less than 100 mSv”. It quoted the US Nuclear Regulatory Commission that “since 1983, the US nuclear industry has monitored more than 100,000 radiation workers each year, and no workers have been exposed to more than 50 mSv in a year since 1989.” A 2012 Massachusetts Institute of Technology study4 exposing mice to low-dose rate radiation for an extended period showed no signs of DNA damage, though a control group receiving the same dose acutely did show damage. This test on live animals confirms other work and epidemiological studies suggesting that people exposed to as much as 1000 mSv/yr at low dose rate will not suffer adverse health effects.
In addition, there is evidence of beneficial effect from low-level radiation (up to about 10 mSv/yr). This ‘radiation hormesis’ may be due to an adaptive response by the body’s cells, the same as that with other toxins at low doses. In the case of carcinogens such as ionizing radiation, the beneficial effect is seen both in lower incidence of cancer and in resistance to the effects of higher doses. This potential hormetic effect is most clearly evident in the data (see Appendix) for over 50,000 survivors of the Hiroshima bomb 1.5 to 3 km from the hypocenter, with dose range 1 to 100 mSv, compared with a large control group.
Further research is underway and a debate continues, as helpfully evident in a June 2016 paper published in Biological Theory, titled Epidemiology Without Biology: False Paradigms, Unfounded Assumptions, and Specious Statistics in Radiation Science. Meanwhile standards for radiation exposure continue to be deliberately conservative.
In the USA, The Low-Dose Radiation Research Act of 2015 calls for an assessment of the current status of US and international low-dose radiation research. It also directs the National Academy of Sciences to “formulate overall scientific goals for the future of low-dose radiation research in the United States” and to develop a long-term research agenda to address those goals. The Act arises from a letter from a group of health physicists who pointed out that the limited understanding of low-dose health risks impairs the nation’s decision-making capabilities, whether in responding to radiological events involving large populations such as the 2011 Fukushima accident or in areas such as the rapid increase in radiation-based medical procedures, the cleanup of radioactive contamination from legacy sites and the expansion of civilian nuclear energy.
Fear of radiation effects
The main effect of low-level radiation arises from fear, not the radiation itself. People who are conditioned to fear any level of ionizing radiation tend to take action to avoid it, and those actions are sometimes much more harmful than any exposure to low-dose radiation could be*. Concerns about low doses of radiation from CT scans and X-rays are not only misguided, but may lead to suffering and deaths from avoided or delayed diagnosis. Also therapeutic benefits of nuclear medicine greatly outweigh any harm that might come from the controlled radiation exposure involved.
* After the Chernobyl accident, some pregnant women in Europe sought abortions without any medical justification, the exposure levels being vastly below those likely to have any effects. Sometimes the fear is promoted by misguided governments, as in Japan where maintaining the evacuation of many people beyond a few weeks has resulted in over 1000 deaths, though exposure levels if people had returned to homes would not be hazardous except possibly in some limited areas, easily defined.
Limiting exposure
Public dose limits for exposure from uranium mining or nuclear plants are usually set at 1 mSv/yr above background.
In most countries the current maximum permissible dose to radiation workers is 20 mSv per year averaged over five years, with a maximum of 50 mSv in any one year. This is over and above background exposure, and excludes medical exposure. The value originates from the International Commission on Radiological Protection (ICRP), and is coupled with the requirement to keep exposure as low as reasonably achievable (ALARA) – taking into account social and economic factors.
Radiation protection at uranium mining operations and in the rest of the nuclear fuel cycle is tightly regulated, and levels of exposure are monitored.
There are four ways in which people are protected from identified radiation sources:
• Limiting time. In occupational situations, dose is reduced by limiting exposure time.
• Distance. The intensity of radiation decreases with distance from its source.
• Shielding. Barriers of lead, concrete or water give good protection from high levels of penetrating radiation such as gamma rays. Intensely radioactive materials are therefore often stored or handled under water, or by remote control in rooms constructed of thick concrete or lined with lead.
• Containment. Highly radioactive materials are confined and kept out of the workplace and environment. Nuclear reactors operate within closed systems with multiple barriers which keep the radioactive materials contained.
UNEP notes: “While the release of radon in underground uranium mines makes a substantial contribution to occupational exposure on the part of the nuclear industry, the annual average effective dose to a worker in the nuclear industry overall has decreased from 4.4 mSv in the 1970s to about 1 mSv today. However, the annual average effective dose to a coal miner is still about 2.4 mSv and for other miners about 3 mSv.” The mining figures are probably for underground situations.
About 23 million workers worldwide are monitored for radiation exposure, and about 10 million of these are exposed to artificial sources, mostly in the medical sector where the annual dose averages 0.5 mSv.
Standards and regulation of radiation exposure
Radiation protection standards are based on the conservative assumption that the risk is directly proportional to the dose, even at the lowest levels, though there is no actual evidence of harm at low levels, below about 100 mSv as short-term dose. To the extent that cell damage is made good within a month (say), chronic dose rates up to 100 mSv per month could also be safe, but the standard assumption, called the ‘linear no-threshold (LNT) hypothesis’, discounts the contribution of any such thresholds and is recommended for practical radiation protection purposes only, such as setting allowable levels of radiation exposure of individuals.
LNT was first accepted by the International Commission on Radiological Protection (ICRP) in 1955, when scientific knowledge of radiation effects was less, and then in 1959 by the United Nations Scientific Committee on the Effects of Atomic Radiation (UNSCEAR) as a philosophical basis for radiological protection at low doses, stating outright that “Linearity has been assumed primarily for purposes of simplicity, and there may or may not be a threshold dose”. (Above 100 mSv acute dose there is some scientific evidence for linearity in dose-effect.) From 1934 to 1955 a tolerance dose limit of 680 mSv/yr was recommended by the ICRP, and no evidence of harm from this – either cancer or genetic – had been documented.
The LNT hypothesis cannot properly be used for predicting the consequences of an actual exposure to low levels of radiation and it has no proper role in low-dose risk assessment. For example, LNT suggests that, if the dose is halved from a high level where effects have been observed, there will be half the effect, and so on. This would be very misleading if applied to a large group of people exposed to trivial levels of radiation and even at levels higher than trivial it could lead to inappropriate actions to avert the doses. At Fukushima following the March 2011 accident, maintaining the evacuation beyond a few days did in fact lead to about 1100 deaths, according to the Japan Reconstruction Agencyi.
Much of the evidence which has led to today’s standards derives from the atomic bomb survivors in 1945, who were exposed to high doses incurred in a very short time. In setting occupational risk estimates, some allowance has been made for the body’s ability to repair damage from small exposures, but for low-level radiation exposure the degree of protection from applying LNT may be misleading. At low levels of radiation exposure the dose-response relationship is unclear due to background radiation levels and natural incidence of cancer. However, the Hiroshima survivor data published in 1958 by UNSCEAR for leukaemia (see Appendix) actually shows a reduction in incidence by a factor of three in the dose range 1 to 100 mSv. The threshold for increased risk here is about 400 mSv. This is very significant in relation to concerns about radiation exposure from contaminated areas after the Chernobyl and Fukushima accidents.
The International Commission on Radiological Protection (ICRP), set up in 1928, is a body of scientific experts and a respected source of guidance on radiation protection, though it is independent and not accountable to governments or the UN. Its recommendations are widely followed by national health authorities, the EU and the IAEA. It retains the LNT hypothesis as a guiding principle.
The International Atomic Energy Agency (IAEA) has published international radiation protection standards since 1962. It is the only UN body with specific statutory responsibilities for radiation protection and safety. Its Safety Fundamentals are applied in basic safety standards and consequent Regulations. However, the UN Scientific Committee on the Effects of Atomic Radiation (UNSCEAR) set up in 1955 is the most authoritative source of information on ionizing radiation and its effects.
In any country, radiation protection standards are set by government authorities, generally in line with recommendations by the ICRP, and coupled with the requirement to keep exposure as low as reasonably achievable (ALARA) – taking into account social and economic factors. The authority of the ICRP comes from the scientific standing of its members and the merit of its recommendations.
The three key points of the ICRP’s recommendations are:
• Justification. No practice should be adopted unless its introduction produces a positive net benefit.
• Optimisation. All exposures should be kept as low as reasonably achievable, economic and social factors being taken into account.
• Limitation. The exposure of individuals should not exceed the limits recommended for the appropriate circumstances.
National radiation protection standards are framed for both Occupational and Public exposure categories.
The ICRP recommends that the maximum permissible dose for occupational exposure should be 20 millisievert per year averaged over five years (i.e. 100 millisievert in 5 years) with a maximum of 50 millisievert in any one year. For public exposure, 1 millisievert per year averaged over five years is the limit. In both categories, the figures are over and above background levels, and exclude medical exposure.j
Post-accident exposure
These low exposure levels are achievable for normal nuclear power and medical activities, but where an accident has resulted in radioactive contamination their application has no net health benefit. There is a big difference between what is desirable in the normal planned operation of any plant, and what is tolerable for dealing with the effects of an accident. Here, restrictive dose limits will limit flexibility in managing the situation and thus their application may increase other health risks, or even result in major adverse health effects, as near Fukushima since March 2011 (see earlier endnote). The objective needs to be to minimize the risks and harm to the individual and population overall, rather than focusing on radiation in isolation.
This is recognised to some extent in the occupational health limits set for cleaning up such situations: the IAEA sets 100 mSv as the allowable short-term dose for emergency workers taking vital remedial actions, and 500 mSv as allowable short-term dose for emergency workers taking life-saving actions. At Fukushima, 250 mSv was set as the allowable short-term dose for workers controlling the disabled reactors during 2011. Following NRA consideration of the Fukushima experience, as well as overseas standards and the science, 250 mSv is now the proposed allowable dose in emergency situations in Japan from April 2016.
But even these levels are low, and there has been no corresponding allowance for neighbouring members of the public – ALARA was the only reference criterion regardless of its collateral effects due to prolonging the evacuation beyond a few days. The death toll and trauma from prolonged evacuation at Fukushima were clearly very much greater than the risks of elevated radiation exposure after the first few days.
This led to the IAEA in May 2013 publishing allowable dose rates for members of the public living normally in affected areas, measured 1m above the contaminated ground. A level of 220 mSv/yr over a full year is “safe for everyone” providing any ingested radioactivity is safe. Shorter term, at 40 times this level, 170 mSv for one week is provisionally safe, and at four times the yearly level – 880 mSv – is provisionally safe for one month.
In March 2011, soon after the Fukushima accident, the ICRP said that it “continues to recommend reference levels of 500 to 1000 mSv to avoid the occurrence of severe deterministic injuries for rescue workers involved in an emergency exposure situation.” For members of the public in such situations, it recommends “reference levels for the highest planned residual dose in the band of 20 to 100 millisieverts (mSv)”, dropping to 1 to 20 mSv/yr when the situation is under control.
When making decisions on evacuations, optimization of all the health risks (not just radiation exposure) is required. Rather than ALARA, a concept of As High As Relatively Safe – AHARS – would be more appropriate in dealing with emergency or existing high exposure situations, based on scientific evidence. This would be similar to the 680 mSv Tolerance Dose allowable to 1950s (assuming gamma & beta radiation), and it would take into account the dose rate. This AHARS approach is supported by Allison (2011) and Cuttler (2012 & 2013) among others, and an AHARS level of 1000 mSv/yr or 100 mSv per month is suggested. This would also mean that most or all of the displaced residents from near the Fukushima plant could return home, without any elevated cancer risk.
Another proposal is for As High as Naturally Existent – AHANE. That is, the highest naturally occurring radiation exposures in the world experienced by a large number of people without discernable negative health effects. At Ramsar in Iran, about 2000 people are exposed to at least 250 mSv/yr with no adverse effect. But in Guarapari, Brazil (pop. 73,000), Kerala, India (pop. 100,000), and Yangjiang, China (pop. 80,000) the average exposures are about 50 mSv/year, 38 mSv/year and 35 mSv/year respectively. In all cases, the residents have life expectancies at least as long as their national peers, and cancer rates slightly lower than fellow countrymen. US Navy shipyard worker data confirm the safety of 50 mSv/yr.
Nuclear fuel cycle radiation exposures
The average annual radiation dose to employees at uranium mines (in addition to natural background) is around 2 mSv (ranging up to 10 mSv). Natural background radiation is about 2 mSv. In most mines, keeping doses to such low levels is achieved with straightforward ventilation techniques coupled with rigorously enforced procedures for hygiene. In some Canadian mines, with very high-grade ore, sophisticated means are employed to limit exposure. (See also information page on Occupational Safety in Uranium Mining.) Occupational doses in the US nuclear energy industry – conversion, enrichment, fuel fabrication and reactor operation – average less than 3 mSv/yr.
Reprocessing plants in Europe and Russia treat used fuel to recover useable uranium and plutonium and separate the highly radioactive wastes. These facilities employ massive shielding to screen gamma radiation in particular. Manual operations are carried by operators behind lead glass using remote handling equipment.
In mixed oxide (MOX) fuel fabrication, little shielding is required, but the whole process is enclosed with access via gloveboxes to eliminate the possibility of alpha contamination from the plutonium. Where people are likely to be working alongside the production line, a 25mm layer of perspex shields neutron radiation from the Pu-240. (In uranium oxide fuel fabrication, no shielding is required.)
Interestingly, due to the substantial amounts of granite in their construction, many public buildings including Australia’s Parliament House and New York Grand Central Station, would have some difficulty in getting a licence to operate if they were nuclear power stations.
Accidental radiation exposure (nuclear and other)
(ARS = Acute radiation syndrome)
The death toll from the Chernobyl accident is about 56, that from misuse of radiotherapy and orphan radioisotope sources is less than 50 – half of those from deliberate therapeutic doses due to faulty equipment or procedures.
Three Mile Island – nuclear power reactor
The March 1979 accident at Three Mile Island nuclear power plant in the USA caused some people near the plant to receive very minor doses of radiation, well under the internationally recommended level. Subsequent scientific studies found no evidence of any harm resulting from that exposure. In 1996, some 2,100 lawsuits claiming adverse health effects from the accident were dismissed for lack of evidence. INES rating 5.
Chernobyl – nuclear power reactor
Immediately after the Chernobyl nuclear power plant disaster in 1986, much larger doses were experienced. Apart from the residents of nearby Pripyat, who were evacuated within two days, some 24,000 people living within 15 km of the plant received an average of 450 mSv before they were evacuated. A total of 5200 PBq of radioactivity (iodine-131 equivalent) was released.
In June 1989, a group of experts from the World Health Organization agreed that an incremental long-term dose of 350 mSv should be the criterion for relocating people affected by the 1986 Chernobyl accident. This was considered a “conservative value which ensured that the risk to health from this exposure was very small compared with other risks over a lifetime”. (For comparison, background radiation averages about 150-200 mSv over a lifetime in most places.)
Out of the 134 severely exposed workers and firemen, 28 of the most heavily exposed died as a result of acute radiation syndrome (ARS) within three months of the accident. Of these, 20 were from the group of 21 that had received over 6.5 Gy, seven (out of 22) had received between 4.2 and 6.4 Gy, and one (out of 50) from the group that had received 2.2-4.1 Gy.5 A further 19 died in 1987-2004 from different causes (see information page on Chernobyl Accident Appendix 2: Health Impacts).
Regarding the emergency workers with doses lower than those causing ARS symptoms, a 2006 World Health Organization report6 referred to studies carried out on 61,000 emergency Russian workers where a total of 4995 deaths from this group were recorded during 1991-1998. “The number of deaths in Russian emergency workers attributable to radiation caused by solid neoplasms and circulatory system diseases can be estimated to be about 116 and 100 cases respectively.” Furthermore, although no increase in leukaemia is discernible yet, “the number of leukaemia cases attributable to radiation in this cohort can be estimated to be about 30.” Thus, 4.6% of the number of deaths in this group are attributable to radiation-induced diseases. (The estimated average external dose for this group was 107 mSv.)
The report also links the accident to an increase in thyroid cancer in children: “During 1992-2000, in Belarus, Russia and Ukraine, about 4000 cases of thyroid cancer were diagnosed in children and adolescents (0-18 years), of which about 3000 occurred in the age group of 0-14 years. For 1152 thyroid cancer patient cases diagnosed among Chernobyl children in Belarus during 1986-2002, the survival rate is 98.8%. Eight patients died due to progression of their thyroid cancer and six children died from other causes. One patient with thyroid cancer died in Russia.”
There has been no increase attributable to Chernobyl in congenital abnormalities, adverse pregnancy outcomes or any other radiation-induced disease in the general population either in the contaminated areas or further afield.
Reports two decades after the accident make it clear that the main health effects from the accident are due to the evacuation of many people coupled with fear engendered, and thousands have died from suicide, depression and alcoholism. The 2006 Chernobyl Forum report said that people in the area suffered a paralysing fatalism due to myths and misperceptions about the threat of radiation, which contributed to a culture of chronic dependency. Some “took on the role of invalids.” Mental health coupled with smoking and alcohol abuse is a very much greater problem than radiation, but worst of all at the time was the underlying level of health and nutrition. Psycho-social effects among those affected by the accident are similar to those arising from other major disasters such as earthquakes, floods and fires.
After the shelterf was built over the destroyed reactor at Chernobyl, a team of about 15 engineers and scientists was set up to investigate the situation inside it. Over several years they repeatedly entered the ruin, accumulating individual doses of up to 15,000 mSv. Daily dose was mostly restricted to 50 mSv, though occasionally it was many times this. None of the men developed any symptoms of radiation sickness, but they must be considered to have a considerably increased cancer risk. INES rating 7.
Fukushima – nuclear power reactors
The March 2011 accident at Fukushima Daiichi nuclear power plant in Japan released about 940 PBq (iodine-131 equivalent) of radioactive material, mostly on days 4 to 6 after the tsunami. In May 2013 UNSCEAR reported that “Radiation exposure following the nuclear accident at Fukushima Daiichi did not cause any immediate health effects. It is unlikely to be able to attribute any health effects in the future among the general public and the vast majority of workers.” The only exception are the 146 emergency workers that received radiation doses of over 100 mSv during the crisis.7 Thyroid doses in children were significantly lower than from the Chernobyl accident. Some 160,000 people were evacuated as a precautionary measure, and prolonging the evacuation resulted in the deaths of about 1100 of them due to stress, and some due to disruption of medical and social welfare facilities. The highest internal radioactivity from ingestion was 12 kBq, some 1000 times less than the level causing adverse health effects at Goiania (see below).
Certainly the main radiation exposure was to workers on site, and the 146 with doses over 100 mSv will be monitored closely for “potential late radiation-related health effects at an individual level.” Six of them had received over 250 mSv – the limit set for emergency workers there, apparently due to inhaling iodine-131 fume early on. There were around 250 workers on site each day. INES rating 7.
Kyshtym, Techa River – military nuclear reprocessing plant
There was a major chemical accident at Mayak Chemical Combine (then known as Chelyabinsk-40) near Kyshtym in Russia in 1957. This plant had been built in haste in the late 1940s for military purposes. The failure of the cooling system for a tank storing many tonnes of dissolved nuclear waste resulted in an ammonium nitrate explosion with a force estimated at about 75 tonnes of TNT (310 GJ). Most of the 740-800 PBq of radioactive contamination settled out nearby and contributed to the pollution of the Techa River, but a plume containing 80 PBq of radionuclides spread hundreds of kilometres northeast. The affected area was already very polluted – the Techa River had previously received about 100 PBq of deliberately dumped waste, and Lake Karachay had received some 4000 PBq. This ‘Kyshtym accident’ killed perhaps 200 people and the radioactive plume affected thousands more as it deposited particularly Cs-127 and Sr-90.
Many people received doses up to 400 mSv at relatively low dose rates from liquid wastes released into the river. This population has shown an increase in cancer rates at levels above 200 mSv. But below this level, cancer incidence falls below the LNT expectations. INES rating 6.
Goiania – orphan source
In 1987 at Goiania8 in Brazil, a discarded radiotherapy source stolen from an abandoned hospital and broken open caused four deaths, 20 cases of radiation sickness and significant contamination of many more. The teletherapy source contained 93 grams of caesium-137 (51 TBq) encased in a shielding canister 51 mm diameter and 48 mm long made of lead and steel, with an iridium window. Various people came in contact with the source over two weeks as it was relayed to a scrapyard, and some were seriously affected. The four deaths (4-5 Sv dose) were family and employees of the scrapyard owner, and 16 others received more than 500 mSv dose. Overall 249 people were found to have significant levels of radioactive material in their bodies. In the 25 years since 1987 there have been zero cancers from radiation among those 249 people affected at Goiania, in spite of the ingestion of up to 100 MBq giving doses as high as 625 mSv/month (8 individuals had higher activity than 100 MBq internally of whom 4 died of Acute Radiation Syndrome but none of cancer). Two healthy babies were born, one to a mother among the most highly contaminated. However fear of the contamination has been the cause of severe stress and depression. In March 2012 Yukiya Amano, the Director General of IAEA, described Goiania as the best illustration of the effect of a terrorist dirty bomb – a few deaths but widespread fear and stress. INES rating 5.
Fleurus – commercial irradiation
In March 2006 at the Institute for Radioelements (IRE) in Fleurus9 in Belgium a worker at a commercial irradiation facility received a high radiation dose of about 4.6 Sv from cobalt-60, resulting in severe health effects. INES rating 5.
In August 2008 about 45 GBq of iodine-131 was released through the stack of the Institute for Radioelements (IRE) in Fleurus, Belgium. The release occurred following the transfer of liquid waste from one tank to another. INES rating 3.
Stamboliysky – commercial irradiation
In June 2011 in Bulgaria, preparations for the recharging of a gamma-irradiation facility with cobalt-60 sources were being undertaken at Stamboliysky. A device already recharged with sources had been taken out, instead of an empty one due to personnel error. As a result, four workers were exposed to a powerful gamma radiation for approximately five minutes, giving them effective doses of over 1 Sv. INES rating 3.
Mexico City – orphan source
In 1962 a boy took home an unshielded cobalt-60 radiography source. Nine people suffered ARS and four of them died from it.
Therac-25 – radiotherapy
In 1985, due to software failure and a fundamental design flaw in the Therac-25 medical irradiation device produced by Atomic Energy of Canada Limited (AECL), delivering 100 times the beta radiation dose intended, there were 3 fatalities from acute radiation syndrome, and three survivors of ARS.
Morocco – orphan source
In March 1984 an iridium-192 source used for industrial radiography was removed from its shielded container and taken home by a worker. Some 11 people suffered from ARS, of whom eight died of ARS.
Zaragoza – radiotherapy
In December 1990 at a clinic in Spain, at least 27 cancer patients were exposed to very high doses from an incorrectly repaired GE electron accelerator. Eleven of them died of ARS.
Costa Rica – radiotherapy
In 1996 at a hospital in San Jose 114 people received an overdose of radiation from cobalt-60 in therapy, about half suffering from ARS, and 13 died of ARS.
Samut Prakan – orphan source
In Thailand in 2000 a 15.7 TBq orphan cobalt-60 source in a scrap metal yard exposed many people. Ten people were hospitalized with ARS, of whom three died.
Mayapuri – orphan source
In India in April 2010 a cobalt-60 source from university equipment in a scrap metal yard exposed many people. Eight people were hospitalized with ARS, of whom one died.
Time perspective
The health effects of exposure both to radiation and to chemical cancer-inducing agents or toxins must be considered in relation to time. There is cause for concern not only about the effects on people presently living, but also about the cumulative effects of actions today over many generations.
Some radioactive materials which reach the environment decay to safe levels within days, weeks or a few years, while others continue their effect for a long time, as do most chemical cancer-inducing agents and toxins. Certainly this is true of the chemical toxicity of heavy metals such as mercury, cadmium and lead. These of course are a natural part of the human environment anyway, like radiation, but maintain their toxicity forever. The essential task for those in government and industry is to prevent excessive amounts of such toxins harming people, now or in the future. Standards are set in the light of research on environmental pathways by which people might ultimately be affected.
Appendix
UNSCEAR Data on Leukaemia Incidence 1950-57 after Exposure at Hiroshima in 1945 line graph
Data points from the left:
i) A control group of 32,963 people over 3 km from the hypocenter. 273 people per million developed leukaemia.
ii) 32,692 people between 2 and 3 km of the hypocenter, with estimated average exposure of about 20 mSv. 92 people per million developed leukaemia.
iii) and iv) 20,113 people between 1.5 and 2 km of the hypocenter, where average doses “were greater than” 500mSv. The left data point (iii) represents calculated radiation exposure for that zone; the right (iv) represents what is thought to be more a more accurate dose, given the cohort’s other radiation-induced symptoms. 398 people per million developed leukaemia.
v) 8810 people between 1 and 1.5 km of the hypocenter, with estimated average exposure of about 5000 mSv. 3746 people per million developed leukaemia.
vi) 1241 survivors less than 1 km from the hypocenter, where over 50,000 were killed. 12,087 people per million developed leukaemia.
The latency period for leukaemia is less than six months. NB this is a log-log graph, and the green line would otherwise be straight.
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Headless CMS
secure CMS?
by Arso and Ilona
5 min
23. Nov 2021.
It doesn’t matter if it's a traditional or headless CMS; website security is one of the most important aspects of running an online presence.
Every business owner should care about website security, and if they don't, the job of every website developer is to emphasize it. Hacked websites lead to data loss, income loss, credibility, and potential lawsuits. Also, that means countless hours of debugging and repair for you, dear developers.
But, how does hacking usually happen?
The favorite way for hackers and bots to gain access is via the login screen. Honestly, if you're trying to get into someone's account, wouldn't that be the first thing to try (even if you are not a hacker)? Password strength also plays a vital role. We can't emphasize that passwords like "12345678" or "password" are not strong enough. 99.9% of accounts get hacked because of the weak password and not using multi-factor authentication (source: Windows Central). Remember, websites and accounts should be fortresses, unreachable to anyone unknown.
Alongside login, outdated websites are an easy target, too. Using an older coupled CMS version that is obsolete means that the security system has not been updated; therefore, it's not protecting that much. One of the weakest spots is third-party add-ons and plugins, and unfortunately, they are popular among non-tech users. So, why were we talking about all those vulnerable targets for attacking a website? Traditional CMSs like WordPress and Joomla, which most people use for building websites, are code and file-heavy. Hence, they have more material vulnerable to cyberattacks.
Believe it or not, not a day goes by without a compromised site, and some data says there are over 30,000 small business website hacks a day. WordPress, Joomla, and Drupal are popular targets for hackers for two simple reasons: their popularity and open-source code. Also, they tend to be used by less technical users, making them an easy target. According to a recent report on website security by ZDNet, most website hacks are related to vulnerabilities in plugins and themes, misconfiguration issues, and a lack of maintenance by web admins, who forgot to update their content management system (CMS). There is only one proper way to protect sensitive data, and that's with the appropriate encryption. For transferring secure data, web applications can use the secure version of HTTP protocol - HTTPS (Hypertext Transfer Protocol Secure) protocol which uses SSL (Secure Sockets Layer) to protect messages transmitted via the network. Secure Data should be written in an encrypted form and remain during transmission to ensure integrity and confidentiality. But securing the data slightly differs with headless CMS versus traditional.
How headless CMS security differs from traditional CMS?
To understand the security difference between these two CMSs, we first have to know how they work. CMS like WordPress have graphical user interfaces that allow content creators and other non-tech users to create and publish styled templates. The content created is stored within a database and displayed to the end-user or reader based on this pre-defined template. In more technical words, the raw data for a blog post is pulled from the MySQL database by Wordpress's PHP application and pushed to the theme. The theme then converts the content into HTML and styles based on its CSS to let the reader consume it. It's clear now how everything in the traditional CMS is packaged together. The frontend and the backend are codependent. Traditional CMS vs. Headless CMS - The Full Comparison Getting all that critical functionality out-of-the-box does translate into... code. Lots and lots of code, lots and lots of files. Which also means lots and lots of potential vulnerabilities to be exploited. Headless CMS is a bit different. Under the pure headless CMS architecture, content is typically delivered through a content distribution network (CDN) and not through a database, as is the case under older CMS versions. Having the frontend and the backend separated keeps the focus on the content creation and storage, with little to no control on the frontend rendition. Unlike a traditional CMS, which consists of backend storage and frontend presentation layer tightly coupled together, headless CMS parts are not codependent; in other words, they are decoupled. Furthermore, the API publishes headless content as read-only. It can also be placed behind one or more layers of code — perhaps an application layer and a security layer — making it even less vulnerable to attack: security tighter, risk of attacks lower. Another popular method of hacking we haven't mentioned in the beginning is through SQL injection. HeadlessCMS combats by running on a server without SQL or even without being connected to SQL. When a developer truly creates a unique decoupled CMS from scratch, like a headless CMS, nothing about your CMS is a known entity.
Is headless CMS that secure?
Ask any headless CMS user or developer who knows about it, and all will say the same - yes.
Here are the security benefits of headless CMS:
It is less susceptible to DDoS attacks.
Headless CMS consists of a backend layer and connects to different front ends using APIs, thus removing the "head."
No database for content, no security threat – simple
There is no CMS code in the creator's environment. Security with flexibility and easy integration and maintenance becomes a piece of cake compared with previous CMSs, and this means the benefit of less anxiety and more productivity.
Fewer updates
There is no need for updating code every time there's a new release of the headless CMS (provided the API is backward compatible), so there is no situation where a slight change in one component may impact the entire system's security and performance.
Web continuity
People often forget to keep website plugins and themes updated, which unethical hackers seek to exploit. Any breach often means compromised continuity of the web page. With decoupled CMS, all temporary issues can usually be resolved in the background with web working.
Future-proofed
Ultimately, the more secure your CMS is, the smoother your adaptation to future demands will be. For example, adding personal user data like emails or other data-sensitive content is more accessible when the headless CMS is secure.
On the other hand, you should always be careful. It's essential that IT teams only shortlist and select headless CMSs with solid track records and security technologies and protocols that protect against cyber attacks. For instance, you should check if the platform provides authentication and authorization features and throttling features to prevent DDoS attacks. Since headless CMS is API-first CMS, they're likely implementing API security best practices by default, but it's a good idea to inquire about these features. Also, think about server security (for example, requiring HTTPS for network communication). For headless CMS to be genuinely more secure than traditional one, API should follow industry standards, and IT teams should ensure the infrastructure uses security best practices. Sometimes, same as with the traditional CMS, the vulnerability appears because of the individual implementation or oversight, not necessarily the chosen software.
Whatever the choice, security is more vital than ever.
In my opinion, losing your online presence is the worst-case scenario for any digital business owner out there. It is not uncommon for an organization to have a situation in which one moment everything is fine, sales are going crazy, and the very next moment your website is unavailable. Puff, gone. Usually, business websites are the target of attacks. In some cases, the attackers take visitors' personal information, causing a crisis. Sometimes they use DDoS attacks to flood a website server with a large amount of traffic and cause it to crash.
Whatever is the case, security must not be neglected in any CMS.
Benjamin Franklin said: "By failing to prepare, you are preparing to fail." Although web security is more vital than ever for today's businesses, with the number of cyberattacks and security breaches skyrocketing, it's not a trial task. It involves a fair amount of planning and executing a complete strategy that goes way beyond simply securing a single or even a dozen websites, plus APIs and development, staging, and production servers. Honest advice for business owners and content creators would be to find experienced developers to make your web or app safer and make your digital experience less stressful. Honest advice for every web developer would be - to give headless cms a chance. The advantage of headless CMSs is the separation of frontend and backend. The reason a decoupled CMS is more secure is that content cannot be accessed by CMS database hackers; it has nothing to exploit. As simple as that.
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Tickling of the Throat: 8 Potential Causes and Home Remedies
What Causes Tickling in my Throat? Causes, Treatment and FAQs
by Saima
Tickling of the Throat
Tickling throat (itching sensation) can be due to an infection in the upper respiratory tract, inflammation in the body, or environmental irritants. It feels like a persistent dusty or prickling sensation that may or may not accompany coughing. The tickly throat may also induce a scratch-like feel at the back of your throat.
Knowing the root cause of the constant itchiness can help eliminate the annoying throat tickle. Let’s dive deep to explore the symptoms, causative agents, treatments, and home remedies for the tickly throat!
What is meant by a Tickling Throat?
The tickly throat is an infectious disorder caused mainly by sore throat, inflammation of the voice box, or sinusitis. It can also indicate a viral and bacterial attack or may result from environmental allergens.
Symptoms of a Tickling Throat
You can feel constant tickly throat alone, or it may come with other symptoms such as:
• The urge-to-cough
• Tickling in the throat with persistent dry cough
• Runny nose
• Hoarseness
• Obstruction in breathing due to nasal congestion
• Inflammation in the upper respiratory tract
• Facial stress
• Fever and chills
• Dripping sensation at the back of the throat
What Causes Tickle in my Throat?
There’re various reasons behind the tickling sensation of the throat. Usually, respiratory tract infections, environmental irritants, or inflammation in the nasal cavities are the most common reasons for the itchy throat. Let’s explore more about its potential causes!
Adult-Onset Allergies
Hypersensitivity to certain irritants in the environment may cause allergies. The hypersensitivity signals the immune system to respond to these irritants in the form of sneezing, itching, tickly throat, runny nose, or trouble breathing.
Dust, pollens, latex, certain foods, or pets are common allergens that can lead to tingly throat and coughing. Approximately, 30-56% of US adults suffer from environmental allergies at some stage of their lives, and about 29% of people have to avoid some foods that they believe trigger their allergy problems.
Sometimes, traveling to a different geographical region can also trigger allergic reactions. Similarly, genetics or family history may also contribute to causing allergies.
Pharyngitis
Usually, inflammation of the throat or pharynx can result in pharyngitis or sore throat. It may also feel like stiff hair in the throat and swelling accompanied by a dry cough. Certain bacterial or viral infections can cause pharyngitis resulting in a dry, scratchy, tickly throat and difficulty swallowing.
Laryngitis
Laryngitis is the inflammation of your voice box or larynx due to viral, bacterial, and fungal infections or environmental pollutants. Hoarseness and painful throat are the most common symptom of laryngitis.
Asthma
Asthma can be a hereditary disease or may be caused by certain environmental pollutants. Chronic infection or inflammation of the airways can cause asthma with symptoms of chest tightness, tickling cough, shortness of breath, and wheezing. About 7.8% of people have asthma in the US, and more common in women and children (aged 5 to 14 years). The highest prevalence of asthma is at the age of 15 to 34 years.
Environmental Irritants
Some environmental allergens such as smoke, dust, pollens, toxins, or chemical substances can trigger the allergic response. These causative agents may cause Type-1 Hypersensitivity involving breathing problems, scratchy throat, tingling of different body parts, inflammation, skin lesions, or lips sensitivity.
Learn more: What Causes Tingling in the Body
Protecting yourself from allergy-triggering factors is the safest way to prevent allergic responses. For this purpose, identify and avoid all the environmental irritants that may cause recurring infections.
Sinusitis and tickling of the Throat
A throat tickle may be caused by sinusitis, pressure on facial regions, nasal congestion, constant headache, and chronic cough. Sinusitis occurs due to viral, bacterial, or fungal infections. Usually, sinus infection issues can appear at any time, but this problem can trigger due to seasonal allergies or environmental toxins.
Coronavirus Disease (COVID-19)
Common symptoms of COVID-19 include itchy throat, trouble breathing, fatigue, fever, dry cough, loss of smell and taste, and severe muscle cramps or pain. However, symptoms of Coronavirus disease may vary from mild to severe depending on the strength of the immune system, age, and health of sufferers.
When you test positive for COVID-19, the first thing to do is quarantine yourself to protect other people in your surroundings from infection. If your symptoms get worse with time, you should consult with a medical emergency as soon as possible.
Gastroesophageal reflux disease (GERD)
GERD is a more serious condition of acid reflux and heartburn. This problem occurs due to the weakening of the esophagus muscles causing acid reflux from the stomach to the esophagus or back side of the throat. Persistent attacks of GERD may severely damage the esophagus line causing bleeding in the esophagus or a burning sensation and pain in the chest area.
How to Stop or Get Rid of Throat Tickle?
Throat tickles can be treated by following some home remedies; however, you must consult your health care advisor in case of severe infections such as COVID-19, high fever with chills, and breathing problems. A few tips are mentioned below for quick ease of this prickling sensation:
• Over-the-counter medications can also help to relieve the pain; however, they might have temporary effect. Non-steroidal, anti-inflammatory medicines can help reduce sore and tingly throat.
• Gargling twice daily with lukewarm water by adding around a quarter teaspoon of salt provides a soothing effect to the inflamed and painful throat area.
• Throat lozenges can help to temporarily relieve the stiff, prickly feeling in the throat because medicated lozenges improve lubrication of your throat by increasing saliva production.
• Plenty of Rest is another excellent way to calm your body while fighting disease. Ensure sleeping 6-8 hours to prevent excessive body stress.
• Herbal Tea and Beverages can help eliminate toxins from your body and proved to be fantastic for relieving dry cough. You can use honey water (one tablespoon of honey in one glass of lukewarm water) for quick relief. Hot ginger tea, turmeric milk, and herbal tea (chamomile, licorice, ginkgo, stinging nettles) are the most effective home remedies for itchy throat and persistent cough.
• Identify the allergens triggering your health problems if you’re suffering from seasonal allergies. You should know the types of foods, environment, medicines, or other household products that worsen your allergy. Preventing these irritants can be much more effective for tickling sensations in the throat and other associated health problems.
Why Won’t Tickle in My Throat go away?
Generally, the causative agents for tickly throat can take a long time (maybe up to 4-6 weeks) for a complete cure. Moreover, if the reason is seasonal allergies, then you may get relief after the end of that particular season or environmental allergen. Moreover, some respiratory infections may also take a few weeks to recover.
Read the article: Causes of Tickling Feet
Persistent Cough with Tickly Throat
Tickling is an infection of the upper respiratory airways caused by specific allergies or diseases. Some coughs are viral, while others may be due to smoking or surrounding pollens. You can suffer from persistent coughing whenever exposed to these irritants.
Is Adult-onset Allergy and Common Cold Same?
Usually, an adult-onset allergy or new-onset seasonal allergies may be confused with the common cold, but they’re quite different. Adult-onset seasonal allergy is not as easy to diagnose and treat as the common cold (can also get better without medication). Moreover, seasonal allergies can’t be treated completely, but you can manage or control their symptoms with care. Nasal sprays, antihistamines, and immunotherapy can significantly reduce the painful symptoms of adult-onset allergy.
Concluding Remarks
Tickling in the throat could have a variety of causes. It could be triggered by a minor issue, such as a cold or an allergen. Sometimes, it might indicate a more severe disease like throat cancer or GERD.
You can try home remedies to stop the tickling. However, you should immediately visit a doctor if the problem persists with more severe symptoms.
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MAN PAGES > SECTION 3
getpwnam_r(3)
getpwnam, getpwnam_r, getpwuid, getpwuid_r - get password file entry
Current Version:
Linux Kernel - 3.80
Synopsis
#include <sys/types.h>
#include <pwd.h>
struct passwd *getpwnam(const char *name);
struct passwd *getpwuid(uid_t uid);
int getpwnam_r(const char *name, struct passwd *pwd,
char *buf, size_t buflen, struct passwd **result); int getpwuid_r(uid_t uid, struct passwd *pwd,
char *buf, size_t buflen, struct passwd **result);
Feature Test Macro Requirements for glibc (see feature_test_macros(7)):
getpwnam_r(), getpwuid_r():
_POSIX_C_SOURCE >= 1 || _XOPEN_SOURCE || _BSD_SOURCE || _SVID_SOURCE || _POSIX_SOURCE
Description
The getpwnam() function returns a pointer to a structure containing the broken-out fields of the record in the password database (e.g., the local password file /etc/passwd, NIS, and LDAP) that matches the username name.
The getpwuid() function returns a pointer to a structure containing the broken-out fields of the record in the password database that matches the user ID uid.
The passwd structure is defined in <pwd.h> as follows:
struct passwd {
char *pw_name; /* username */
char *pw_passwd; /* user password */
uid_t pw_uid; /* user ID */
gid_t pw_gid; /* group ID */
char *pw_gecos; /* user information */
char *pw_dir; /* home directory */
char *pw_shell; /* shell program */
};
See passwd(5) for more information about these fields.
The getpwnam_r() and getpwuid_r() functions obtain the same information as getpwnam() and getpwuid(), but store the retrieved passwd structure in the space pointed to by pwd. The string fields pointed to by the members of the passwd structure are stored in the buffer buf of size buflen. A pointer to the result (in case of success) or NULL (in case no entry was found or an error occurred) is stored in *result.
The call
sysconf(_SC_GETPW_R_SIZE_MAX)
returns either -1, without changing errno, or an initial suggested size for buf. (If this size is too small, the call fails with ERANGE, in which case the caller can retry with a larger buffer.)
Return Value
The getpwnam() and getpwuid() functions return a pointer to a passwd structure, or NULL if the matching entry is not found or an error occurs. If an error occurs, errno is set appropriately. If one wants to check errno after the call, it should be set to zero before the call.
The return value may point to a static area, and may be overwritten by subsequent calls to getpwent(3), getpwnam(), or getpwuid(). (Do not pass the returned pointer to free(3).)
On success, getpwnam_r() and getpwuid_r() return zero, and set *result to pwd. If no matching password record was found, these functions return 0 and store NULL in *result. In case of error, an error number is returned, and NULL is stored in *result.
Errors
0 or ENOENT or ESRCH or EBADF or EPERM or ...
The given name or uid was not found.
EINTR
A signal was caught.
EIO
I/O error.
EMFILE
The maximum number (OPEN_MAX) of files was open already in the calling process.
ENFILE
The maximum number of files was open already in the system.
ENOMEM
Insufficient memory to allocate passwd structure.
ERANGE
Insufficient buffer space supplied.
Files
/etc/passwd
local password database file
Attributes
Multithreading (see pthreads(7))
The getpwnam() and getpwuid() functions are not thread-safe.
The getpwnam_r() and getpwuid_r() functions are thread-safe.
Conforming To
SVr4, 4.3BSD, POSIX.1-2001. The pw_gecos field is not specified in POSIX, but is present on most implementations.
Notes
The formulation given above under "RETURN VALUE" is from POSIX.1-2001. It does not call "not found" an error, and hence does not specify what value errno might have in this situation. But that makes it impossible to recognize errors. One might argue that according to POSIX errno should be left unchanged if an entry is not found. Experiments on various UNIX-like systems show that lots of different values occur in this situation: 0, ENOENT, EBADF, ESRCH, EWOULDBLOCK, EPERM, and probably others.
The pw_dir field contains the name of the initial working directory of the user. Login programs use the value of this field to initialize the HOME environment variable for the login shell. An application that wants to determine its user's home directory should inspect the value of HOME (rather than the value getpwuid(getuid())->pw_dir) since this allows the user to modify their notion of "the home directory" during a login session. To determine the (initial) home directory of another user, it is necessary to use getpwnam(username)->pw_dir or similar.
Example
The program below demonstrates the use of getpwnam_r() to find the full username and user ID for the username supplied as a command-line argument.
#include <pwd.h>
#include <stdio.h>
#include <stdlib.h>
#include <unistd.h>
#include <errno.h>
int
main(int argc, char *argv[])
{
struct passwd pwd;
struct passwd *result;
char *buf;
size_t bufsize;
int s;
if (argc != 2) {
fprintf(stderr, "Usage: %s username\n", argv[0]);
exit(EXIT_FAILURE);
}
bufsize = sysconf(_SC_GETPW_R_SIZE_MAX);
if (bufsize == -1) /* Value was indeterminate */
bufsize = 16384; /* Should be more than enough */
buf = malloc(bufsize);
if (buf == NULL) {
perror("malloc");
exit(EXIT_FAILURE);
}
s = getpwnam_r(argv[1], &pwd, buf, bufsize, &result);
if (result == NULL) {
if (s == 0)
printf("Not found\n");
else {
errno = s;
perror("getpwnam_r");
}
exit(EXIT_FAILURE);
}
printf("Name: %s; UID: %ld\n", pwd.pw_gecos, (long) pwd.pw_uid);
exit(EXIT_SUCCESS);
}
See Also
Colophon
This page is part of release 3.80 of the Linux man-pages project. A description of the project, information about reporting bugs, and the latest version of this page, can be found at http://www.kernel.org/doc/man-pages/.
License & Copyright
Copyright 1993 David Metcalfe ([email protected])
and Copyright 2008, Linux Foundation, written by Michael Kerrisk
%%%LICENSE_START(VERBATIM)
Permission is granted to make and distribute verbatim copies of this
manual provided the copyright notice and this permission notice are
preserved on all copies.
Permission is granted to copy and distribute modified versions of this
manual under the conditions for verbatim copying, provided that the
entire resulting derived work is distributed under the terms of a
permission notice identical to this one.
Since the Linux kernel and libraries are constantly changing, this
manual page may be incorrect or out-of-date. The author(s) assume no
responsibility for errors or omissions, or for damages resulting from
the use of the information contained herein. The author(s) may not
have taken the same level of care in the production of this manual,
which is licensed free of charge, as they might when working
professionally.
Formatted or processed versions of this manual, if unaccompanied by
the source, must acknowledge the copyright and authors of this work.
%%%LICENSE_END
References consulted:
Linux libc source code
Lewine's "POSIX Programmer's Guide" (O'Reilly & Associates, 1991)
386BSD man pages
Modified 1993-07-24 by Rik Faith ([email protected])
Modified 1996-05-27 by Martin Schulze ([email protected])
Modified 2003-11-15 by aeb
2008-11-07, mtk, Added an example program for getpwnam_r().
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how to model an image of landscape appropriately?
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is there a way to 3d plot a 2d image like from this /matlabcentral/answers/uploaded_files/6759/7.jpg
to this
/matlabcentral/answers/uploaded_files/6760/2.jpg and ploting should rely on image data. thanks
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How to Boot Virtual Machine in the Recovery Mode
Many are the ways through which you can boot a virtual machine in the recovery mode. Here, we are going to focus on of the most effective ways of achieving the same. We will inform you of a method that geeks and other computer experts use to go around the task. Every step below is critical therefore; there is need for you to pay close attention to each one of them.
What are the steps to take?
In this method, the buttons from F1, F2… etc. are very important so there is need to ensure that all of them are functional before beginning the process. Below are the steps to take:
1. 1. Opening of the OS X System preferences
The first step involves opening of the OS X system preferences using your keyboard. Enable all the keys that will serve as standard functional keys.
1. 2. Open the configuration pane
The next step involves opening of the virtual machine configuration pane. You should follow the steps from hardware to boot order where you will get various options. Click on the “select boot device on startup” option. You should then start your virtual machine and make its window active be clicking on it.
1. 3. Press F8
Depending on the type of operating system that you are using, press the F8 button before anything appears on the screen your computer. For some types of OSs like Windows 8.1, you will have to press Shift and F8 to achieve the same effect as the other types.
It is important to note that this stage requires patience. We recommend that after booting your virtual machine; continue pressing on F8 multiple times until the advanced boot option appears on your screen. If you do not achieve this on the first attempt, shut down the computer and redo the procedure.
1. 4. Choose repair your computer option
When the advanced boot option appears on your screen, use the arrow keys to select the “repair your computer option”. Windows will then ask you to choose your language and from here, you can enter your credentials and continue. If you are using Windows 8.1, choose the troubleshoot option.
1. 5. Select boot device on startup option
After the repair or the troubleshooting process is over, the next step is to select the boot device on startup option.It is important to note that you might not be able to boot your machine into recovery mode if you do not have a system recovery partition. However, you can go about that by booting the virtual machine from the installation media. Booting a virtual machine in the recovery mode is not as hard a task. All you need to do is follow the steps above.
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RNA
Like DNA, RNA (ribonucleic acid) is essential for all known forms of life. RNA monomers are also nucleotides. Unlike DNA, RNA in biological cells is predominantly a single-stranded molecule. While DNA contains deoxyribose, RNA contains ribose, characterised by the presence of the 2'-hydroxyl group on the pentose ring (Figure 5). This hydroxyl group make RNA less stable than DNA because it is more susceptible to hydrolysis. RNA contains the unmethylated form of the base thymine called uracil (U) (Figure 6), which gives the nucleotide uridine.
Chemical structure of RNA
Figure 5 Chemical structure of RNA: nucleotides containing a ribose sugar (carbons numbered 1' through 5') with a base attached to the 1' position. From wikipedia.
Nucleobase Chemical structure of adenine
Adenine
Chemical structure of guanine
Guanine
Chemical structure of thymine
Thymine
Chemical structure of cytosine
Cytosine
Chemical structure of uracil
Uracil
Figure 6: Nucleobase structure of RNA and DNA.
RNA performs a variety of functions in the cell. Messenger RNA (mRNA) carries the genetic information that directs the synthesis of proteins. Some viruses use RNA instead of DNA as their genetic material. Most of the RNA, however, does not code for proteins. These RNAs are called noncoding (ncRNA) and can be encoded by their own RNA genes or can derive from mRNA intronsTransfer RNA (tRNA) and ribosomal RNA (rRNA) are involved in the process of translation. There are also non-coding RNAs involved in gene regulation, RNA processing and other processes.
Most RNA molecules contain short self-complementary sequences that fold and pair with each other into highly structured forms. These base-pairing interactions are part of RNA secondary structure. The unpaired regions form structures such as hairpin loops, bulges and internal loops, which may be of functional importance (Figure 7). Examples include Rho-independent terminator stem-loops and the tRNA cloverleaf.
Secondary and tertiary structure of tRNA
Figure 7 Secondary and tertiary structure of tRNA. Unpaired regions are in grey and paired regions are in colour. Image from wikipedia.
The functional form of single-stranded RNA molecules, just like proteins, frequently requires a specific tertiary (3D) structure.RNA can also form RNA-RNA and DNA-RNA duplexes. Most RNA structures in the Protein Data Bank (PDB) (archive of macromolecular structural data) (3) contain double-stranded RNA folded into tertiary structures.
Some RNA structures provide binding sites for other molecules and have chemically active centres. An example, (Figure 8) is the molecular recognition of vitamin B12 by an RNA structure (4). Vitamin B12 binding to RNA regulates hepatitis C virus function (5).
The structure of vitamin B12 bound to RNA
Figure 8 The structure of vitamin B12 bound to RNA. The molecular recognition is achieved by the folding of an initially unstructured RNA around its ligand.
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=pod =head1 NAME SQL::Amazon::UserGuide - User Guide for DBD/SQL::Amazon =head1 SQL::Amazon SYNOPSIS # # create the parser, passing in the current Amazon metadata # my $parser = SQL::Amazon::Parser->new(\%attrs); # # parse a SQL statement, returning a SQL::Amazon::Statement # my $stmt = $parser->parse($sql_stmt) or die "Parse failed: " . $parser->errstr; # # evaluate the parse tree, using an evaluation object # for driver specific evaluation # returns either a scalar rowcount for write operations, # or a SQL::Amazon::Table object for SELECT # or undef on error # my $results = $stmt->execute($params) or die "Evaluation failed: " . $stmt->errstr; =head1 DESCRIPTION SQL::Amazon provides the various components required by B L to query the Amazon E-Commerce Service 4.0 I using SQL. B. The suite includes the following components: =over 4 =item B provides SQL parsing and query plan generation. Implemented as a subclass of L, part of the L bundle. =item B provides SQL query plan execution. Implemented as a subclass of L. =item B provides SQL::Amazon-specific predicate functions, including B. =item B provides a factory class for generating SQL::Amazon::Request::Request objects based on the predicates in a query's WHERE clause. =item B provides a temporary storage object for intermediate results extracted from the base table cache objects. Acts as a L object for L processing. =item B provides a global storage engine for managing data caching and retrieval. =item B provides a base class for all ECS request objects, including numerous default method implementations for building and sending requests, and processing the responses into the base table cache objects. =item B a subclass of SQL::Amazon::Request::Request for the ItemLookup request; also acts as a base class for the ItemSearch request. =item B a subclass of SQL::Amazon::Request::ItemLookup for the ItemSearch request =item B provides a base class for table cache objects, including methods for data type conversion, keyed lookup, and cache management. =item B> provides table-specific implementations of the Table base class. =back =head2 Prerequisites Perl 5.8.0 L 1.10 L 0.15 L 1.47 L 2.09 L I> =head2 SQL Dialect SQL::Amazon::Parser currently supports the same subset of standard SQL as supported by L, using the L dialect, with the following additional predicate functions for keyword searches: =over 4 =item BEIEB<)> Results in generation of an ECS ItemSearch request, with either the Keywords parameter (for MATCHES ALL), Power parameter (for MATCHES ANY), or TextStream parameter (for MATCHES TEXT). =item BEIEB<)> Results in generation of an ECS ItemSearch request, with either Power parameters. Individual items in the I are wrapped in parentheses, and joined with 'and' operators. =back =head2 Special SQL Considerations =head4 Mapping ECS Hierarchical Response Data to Relational Structures The hierarchical structure of the data returned by the Amazon ECS requires some manipulation by SQL::Amazon to conform to a flattened relational model. As a result, queries against the B table may cause multiple tables to be populated, e.g., B, etc. This behavior is especially true when the B parameter is set to B I<(the default)>, in which case all detail information about matching books is returned. As a result, some queries require the use of JOINs to properly reflect the data model abstraction. E.g., When searching for only New books, the B ECS request parameter must be applied; however, B is a column in the B table, B the B table. Hence, a JOIN between the B table and B table on the B column is required, along with a predicate test for B. The JOIN is required, even though the returned B
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从JS对象开始，谈一谈“不可变数据”和函数式编程
作为前端开发者，你会感受到JS中对象(Object)这个概念的强大。我们说“JS中一切皆对象”。最核心的特性，例如从String，到数组，再到浏览器的APIs，“对象”这个概念无处不在。在这里你可以了解到JS Objects中的一切。
同时，随着React的强势崛起，不管你有没有关注过这个框架，也一定听说过一个概念—不可变数据(immutable.js)。究竟什么是不可变数据？这篇文章会从JS源头—对象谈起，让你逐渐了解这个函数式编程里的重要概念。
JS中的对象是那么美妙：我们可以随意复制他们，改变并删除他们的某项属性等。但是要记住一句话：
“伴随着特权，随之而来的是更大的责任。”
(With great power comes great responsibility)
的确，JS Objects里概念太多了，我们切不可随意使用对象。下面，我就从基本对象说起，聊一聊不可变数据和JS的一切。
这篇文章缘起于Daniel Leite在2017年3月16日的新鲜出炉文章：Things you should know about Objects and Immutability in JavaScript，我进行了大致翻译并进行大范围“改造”，同时改写了用到的例子，进行了大量更多的扩展。
“可变和共享”是万恶之源
不可变数据其实是函数式编程相关的重要概念。相对的，函数式编程中认为可变性是万恶之源。但是，为什么会有这样的结论呢？
这个问题可能很多程序员都会有。其实，如果你的代码逻辑可变，不要惊慌，这并不是“政治错误”的。比如JS中的数组操作，很对都会对原数组进行直接改变，这当然并没有什么问题。比如：
let arr = [1, 2, 3, 4, 5];
arr.splice(1, 1); // 返回[2];
console.log(arr); // [1, 3, 4, 5];
这是我们常用的“删除数组某一项”的操作。好吧，他一点问题也没有。
问题其实出现在“滥用”可变性上，这样会给你的程序带来“副作用”。先不必关心什么是“副作用”，他又是一个函数式编程的概念。
我们先来看一下代码实例：
const student1 = {
school: 'Baidu',
name: 'HOU Ce',
birthdate: '1995-12-15',
}
const changeStudent = (student, newName, newBday) => {
const newStudent = student;
newStudent.name = newName;
newStudent.birthdate = newBday;
return newStudent;
}
const student2 = changeStudent(student1, 'YAN Haijing', '1990-11-10');
// both students will have the name properties
console.log(student1, student2);
// Object {school: "Baidu", name: "YAN Haijing", birthdate: "1990-11-10"}
// Object {school: "Baidu", name: "YAN Haijing", birthdate: "1990-11-10"}
我们发现，尽管创建了一个新的对象student2，但是老的对象student1也被改动了。这是因为JS对象中的赋值是“引用赋值”，即在赋值过程中，传递的是在内存中的引用(memory reference)。具体说就是“栈存储”和“堆存储”的问题。具体图我就不画了，理解不了可以单找我。
不可变数据的强大和实现
我们说的“不可变”，其实是指保持一个对象状态不变。这样做的好处是使得开发更加简单，可回溯，测试友好，减少了任何可能的副作用。
函数式编程认为：
只有纯的没有副作用的函数，才是合格的函数。
好吧，现在开始解释下“副作用”(Side effect)：
在计算机科学中，函数副作用指当调用函数时，除了返回函数值之外，还对主调用函数产生附加的影响。例如修改全局变量（函数外的变量）或修改参数。
－维基百科
函数副作用会给程序设计带来不必要的麻烦，给程序带来十分难以查找的错误，并降低程序的可读性。严格的函数式语言要求函数必须无副作用。
那么我们避免副作用，创建不可变数据的主要实现思路就是：一次更新过程中，不应该改变原有对象，只需要新创建一个对象用来承载新的数据状态。
我们使用纯函数(pure functions)来实现不可变性。纯函数指无副作用的函数。
那么，具体怎么构造一个纯函数呢？我们可以看一下代码实现，我对上例进行改造：
const student1 = {
school: "Baidu",
name: 'HOU Ce',
birthdate: '1995-12-15',
}
const changeStudent = (student, newName, newBday) => {
return {
...student, // 使用解构
name: newName, // 覆盖name属性
birthdate: newBday // 覆盖birthdate属性
}
}
const student2 = changeStudent(student1, 'YAN Haijing', '1990-11-10');
// both students will have the name properties
console.log(student1, student2);
// Object {school: "Baidu", name: "HOU Ce", birthdate: "1995-12-15"}
// Object {school: "Baidu", name: "YAN Haijing", birthdate: "1990-11-10"}
需要注意的是，我使用了ES6中的解构(destructuring)赋值。
这样，我们达到了想要的效果：根据参数，产生了一个新对象，并正确赋值，最重要的就是并没有改变原对象。
创建纯函数，过滤副作用
现在，我们知道了“不可变”到底指的是什么。接下来，我们就要分析一下纯函数应该如何实现，进而生产不可变数据。
其实创建不可变数据方式有很多，在使用原生JS的基础上，我推荐的方法是使用现有的Objects API和ES6当中的解构赋值（上例已经演示）。现在看一下Objects.assign的实现方式：
const student1 = {
school: "Baidu",
name: 'HOU Ce',
birthdate: '1995-12-15',
}
const changeStudent = (student, newName, newBday) => Object.assign({}, student, {name: newName, birthdate: newBday})
const student2 = changeStudent(student1, 'YAN Haijing', '1990-11-10');
// both students will have the name properties
console.log(student1, student2);
// Object {school: "Baidu", name: "HOU Ce", birthdate: "1995-12-15"};
// Object {school: "Baidu", name: "YAN Haijing", birthdate: "1990-11-10"};
同样，如果是处理数组相关的内容，我们可以使用：.map, .filter或者.reduce去达成目标。这些APIs的共同特点就是不会改变原数组，而是产生并返回一个新数组。这和纯函数的思想不谋而合。
但是，再说回来，使用Object.assign请务必注意以下几点：
1）他的复制，是将所有可枚举属性，复制到目标对象。换句话说，不可枚举属性是无法完成复制的。
2）对象中如果包含undefined和null类型内容，会报错。
3）最重要的一点：Object.assign方法实行的是浅拷贝，而不是深拷贝。
第三点很重要，也就是说，如果源对象某个属性的值是对象，那么目标对象拷贝得到的是这个属性对象的引用。这也就意味着，当对象存在嵌套时，还是有问题的。比如下面代码：
const student1 = {
school: "Baidu",
name: 'HOU Ce',
birthdate: '1995-12-15',
friends: {
friend1: 'ZHAO Wenlin',
friend2: 'CHENG Wen'
}
}
const changeStudent = (student, newName, newBday, friends) => Object.assign({}, student, {name: newName, birthdate: newBday})
const student2 = changeStudent(student1, 'YAN Haijing', '1990-11-10');
// both students will have the name properties
console.log(student1, student2);
// Object {school: "Baidu", name: "HOU Ce", birthdate: "1995-12-15", friends: Object}
// Object {school: "Baidu", name: "YAN Haijing", birthdate: "1990-11-10", friends: Object}
student2.friends.friend1 = 'MA xiao';
console.log(student1.friends.friend1); // "MA xiao"
对student2 friends列表当中的friend1的修改，同时也影响了student1 friends列表当中的friend1。
JS本身的苍白无力VS不可变数据类库
以上，我们分析了纯JS如何实现不可变数据。这样处理带来的一个负面影响在于：一些经典APIs都是shallow处理，比如上文提到的Object.assign就是典型的浅拷贝。如果遇到嵌套很深的结构，我们就需要手动递归。这样做呢，又会存在性能上的问题。
比如我自己动手用递归实现一个深拷贝，需要考虑循环引用的“死环”问题，另外，当使用大规模数据结构时，性能劣势尽显无疑。我们熟悉的jquery extends方法，某一版本（最新版本情况我不太了解）的实现是进行了三层拷贝，也没有达到完备的deep copy。
总之，实现不可变数据，我们必然要关心性能问题。针对于此，我推荐一款已经“大名鼎鼎”的——immutable.js类库来处理不可变数据。
他的实现既保证了不可变性，又保证了性能大限度优化。原理很有意思，下面这段话，我摘自camsong前辈的文章
Immutable实现的原理是Persistent Data Structure（持久化数据结构），也就是使用旧数据创建新数据时，要保证旧数据同时可用且不变。
同时为了避免deepCopy把所有节点都复制一遍带来的性能损耗，Immutable使用了Structural Sharing（结构共享），即如果对象树中一个节点发生变化，只修改这个节点和受它影响的父节点，其它节点则进行共享。
感兴趣的读者可以深入研究下，这是很有意思的。如果有需要，我也愿意再写一篇immutable.js源码分析。
总结
我们使用JavaScript操纵对象，这样的方式很简单便捷。但是，这样操控的基础是在JavaScript灵活机制的熟练掌握上。不然很容易使你“头大”。
在我开发的百度某部门私信项目中，因为使用了React+Redux技术栈，并且数据结构较为负责，所以我也采用了immutable.js实现。
最后，在前端开发中，函数式编程越来越热，并且在某种程度上已经取代了“过程式”编程和面向对象思想。
我的感想是在某些特定的场景下，不要畏惧变化，拥抱未来。
就像我很喜欢的葡萄牙诗人安德拉德一首诗中那样说的：
我同样不知道什么是海，
赤脚站在沙滩上，
急切地等待着黎明的到来。
Happy Coding!
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How much does an alcoholic drink in a day?
How many drinks a day is considered an alcoholic?
Heavy Alcohol Use:
NIAAA defines heavy drinking as follows: For men, consuming more than 4 drinks on any day or more than 14 drinks per week. For women, consuming more than 3 drinks on any day or more than 7 drinks per week.
How much alcohol makes you an alcoholic?
For men, this typically occurs after 5 drinks and 4 drinks for women within a 2-hour time frame. The Substance Abuse and Mental Health Services Administration (SAMHSA) defines heavy alcohol use as 5 or more days in the past month of binge drinking.
What happens if I drink every night?
Daily alcohol use can cause fibrosis or scarring of the liver tissue. It can also cause alcoholic hepatitis, which is an inflammation of the liver. With long-term alcohol abuse, these conditions occur together and can eventually lead to liver failure.
Is drinking 3 beers a day an alcoholic?
According to the National Institute on Alcohol Abuse and Alcoholism, drinking is considered to be in the moderate or low-risk range for women at no more than three drinks in any one day and no more than seven drinks per week. For men, it is no more than four drinks a day and no more than 14 drinks per week.
IT IS IMPORTANT: Best answer: Is anything added to wine?
How do you tell if you’re an alcoholic?
What are the signs or symptoms of alcoholism?
1. A lack of interest in previously normal activities.
2. Appearing intoxicated more regularly.
3. Needing to drink more in order to achieve the same effects.
4. Appearing tired, unwell or irritable.
5. An inability to say no to alcohol.
6. Anxiety, depression or other mental health problems.
What are the first signs of liver damage from alcohol?
Generally, symptoms of alcoholic liver disease include abdominal pain and tenderness, dry mouth and increased thirst, fatigue, jaundice (which is yellowing of the skin), loss of appetite, and nausea. Your skin may look abnormally dark or light. Your feet or hands may look red.
Is it normal to drink everyday?
For some individuals, one drink a day may be too much. For others it could be 2 to 3 drinks a day. The Center for Disease Control and Prevention (CDC) recommends up to one drink per day for women and two drinks per day for men. Anything exceeding that could be considered unhealthy.
Is drinking 12 beers a day bad?
In summary, if you’re wondering how many beers a day is safe, the answer for most people is one to two. Drinking more than that on a regular basis can put you at risk, and often reverse any health benefits of drinking beer. It’s a fine line to walk.
Can I drink 6 beers a day?
Large volumes of alcohol — beer, wine or whiskey — can damage the heart, raise blood pressure, wreak havoc on the digestive tract and injure the pancreas. … A man who drinks six to eight 12-ounce cans of beer every day on a regular basis can almost count on developing liver cirrhosis within 10 to 15 years.
IT IS IMPORTANT: You asked: How much alcohol can you transport in Alabama?
Is someone who drinks everyday an alcoholic?
Myth: I don’t drink every day OR I only drink wine or beer, so I can’t be an alcoholic. Fact: Alcoholism is NOT defined by what you drink, when you drink it, or even how much you drink. It’s the EFFECTS of your drinking that define a problem.
Is drinking every night bad?
Having a drink every night isn’t necessarily a bad thing. But, at any level of drinking, be it moderate drinking or heavy alcohol dependence, it’s a smart move to know the risks and stay in control.
How long can you live as an alcoholic?
People hospitalized with alcohol use disorder have an average life expectancy of 47–53 years (men) and 50–58 years (women) and die 24–28 years earlier than people in the general population.
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How Are the Three Subphases of Interphase Alike? What Key Event Happens During the S Phase?
Have you ever wondered what happens to a cell before it divides? That process is known as interphase and it’s divided into three main subphases: G1, S, and G2. These subphases have similar characteristics and are vital for the proper cell division to occur. In this article, we will explore how these three subphases of interphase are alike and what crucial event happens during the S phase.
Firstly, the G1 phase is the initial subphase of interphase characterized by cell growth and preparation for DNA replication. During this phase, the cell undergoes a lot of growth to ensure that it has enough energy and nutrients to duplicate its genetic material. This subphase prepares the cell for the challenging S phase that follows.
The S phase is the focal point of the interphase, and one of the most important events occur here. During this subphase, the cell’s DNA is replicated to produce an exact copy of the genetic material. DNA is the blueprint of life, and each cell needs a copy of the same genetic information to divide and create two daughter cells. Hence, the S phase is crucial for the success of cell division and its proper functioning.
Overall, the subphases of interphase have a lot in common, and they are crucial for cell division to occur successfully. Even though each subphase has its unique characteristics, they depend on one another to create two identical daughter cells at the end of the process. In the next section, we will explore in detail each subphase of the interphase.
Introduction to Interphase
Interphase is the part of the cell cycle where a cell grows, replicates its DNA, and prepares to undergo mitosis or meiosis. It can be divided into three subphases, G1, S, and G2. These three subphases all share some similarities but also have distinct differences. In this article, we will focus on how the three subphases of interphase are alike and what key event happens during the S phase.
Similarities between the three subphases of interphase
• All three subphases involve cell growth and preparation for division
• During all three subphases, the cell prepares for DNA replication
• Protein synthesis occurs during all three subphases
• Metabolic activity is high during all three subphases
The key event during the S phase
The S phase, or synthesis phase, is the period during interphase when DNA replication occurs. During this phase, the DNA strands unwind and separate, and new complementary strands are synthesized. This process ensures that each of the resulting daughter cells will have an identical copy of the genetic material. Without DNA duplication during S phase, a cell could not undergo mitosis or meiosis and divide properly.
Subphase Main events
G1 Cell growth and preparation for DNA replication
S DNA replication
G2 Cell growth and preparation for cell division
In conclusion, the three subphases of interphase have similarities in terms of cell growth, preparation for division, DNA replication, protein synthesis, and metabolic activity. However, the S phase stands out as the key event during interphase, where DNA replication takes place, and no further cell division can occur without it.
Overview of Cell Division
Cell division is the process by which cells divide and multiply. It involves a series of events that take place in a specific order. The process is divided into two main stages: interphase and mitosis.
Interphase is the period between two cell divisions. It is divided into three subphases: G1, S, and G2. During interphase, the cell grows, replicates its DNA, and prepares for mitosis. Mitosis is the process by which the cell divides its nucleus and produces two identical daughter cells.
The Three Subphases of Interphase
• G1: During this subphase, the cell is actively growing and carrying out its normal functions. It prepares for DNA replication by synthesizing various enzymes and proteins needed for the process.
• S: In this subphase, DNA replication takes place. The cell synthesizes new DNA strands and forms new chromosomes. The cell’s nucleus doubles in size, and by the end of this subphase, there are two identical sets of chromosomes.
• G2: During this subphase, the cell continues to grow and prepare for mitosis. It checks for DNA damage and repairs it if necessary. The cell synthesizes proteins needed for cell division, such as microtubules that form the spindle apparatus.
The Key Event During the S Phase
The S phase is the most critical subphase of interphase as it is when DNA replication occurs. The cell synthesizes new DNA strands and forms new chromosomes. The process is highly regulated to ensure that the DNA is accurately replicated, with no errors or mutations. The cell also checks for DNA damage and repairs it if necessary.
Event Description
Initiation Cell initiates DNA replication by unwinding the double helix structure of the DNA molecule.
Elongation The cell synthesizes new DNA strands, forming new chromosomes.
Termination The DNA replication process is terminated, and the two identical sets of chromosomes are formed.
Errors during the S phase can result in mutations and chromosomal abnormalities, leading to various diseases, including cancer. Therefore, the S phase is tightly regulated to ensure that DNA replication occurs correctly.
G1 Phase of Interphase
G1 phase, also known as Gap 1 phase, is the first subphase of interphase and occurs immediately after mitosis. During this phase, the cell grows in size and makes new organelles necessary for cell division. The key event that happens during G1 phase is the replication of mitochondria and the synthesis of ribosomal RNA.
How are the three subphases of interphase alike?
• All three subphases occur before mitosis.
• All three subphases involve the growth and preparation of the cell for division.
• All three subphases are critical for successful cell division.
G2 Phase of Interphase
G2 phase, also known as Gap 2 phase, is the third subphase of interphase that occurs after the S phase. During this phase, the cell continues to grow and prepares for mitosis. The key event that happens during G2 phase is the synthesis of microtubules, which are necessary for cell division.
In addition to this, the cell also double-checks its DNA for any errors that may have occurred during DNA replication. If any errors are found, they are corrected before moving on to mitosis.
The S Phase of Interphase
The S phase, also known as the synthesis phase, is the second subphase of interphase. During this phase, the DNA is replicated to ensure that the new cells that are formed during cell division have identical copies of the genetic material. The key event that happens during S phase is the replication of DNA.
Protein Involved Function
Helicase Unwinds the DNA double helix
Single-stranded binding proteins Stabilizes single-stranded DNA
Primase Synthesizes RNA primers
DNA polymerase III Elongates new DNA strands
DNA polymerase I Removes RNA primers and fills gaps with DNA
Ligase Joins Okazaki fragments on the lagging strand
The replication of DNA ensures that each daughter cell will have a complete set of genetic material. The process involves multiple proteins working together, including helicase, primase, DNA polymerase III, and ligase.
S Phase of Interphase
The S phase of interphase is the second subphase of interphase, which is the period between cell divisions. During this phase, the cells engage in DNA synthesis, which is the process that duplicates the genetic material of a cell. The S phase of interphase is essential for cells to prepare for cell division, as without proper duplication of DNA, the cell cannot divide properly.
• The S phase of interphase is similar to the G1 and G2 subphases, as they all involve cell growth and preparation for cell division.
• All three subphases are critical for the proper functioning of cells, as they prepare the cells for the subsequent process of cell division.
• The S phase is also similar to the M phase of the cell cycle, as both involve genetic material duplication, but the M phase only occurs during cell division.
During the S phase of interphase, DNA replication occurs through a process known as semi-conservative replication. This involves the separation of the DNA strands and the formation of new complementary strands to create a new double-stranded DNA molecule. The S phase is also characterized by an increased rate of protein synthesis, which is necessary to support the process of DNA replication.
A key event during the S phase of interphase is the replication of chromosomes, which is crucial for cell division. Chromosome replication leads to the formation of two identical sister chromatids, which contain the same genetic information. These sister chromatids are joined at a structure known as the centromere, and they will separate during cell division to ensure that each daughter cell receives an identical copy of the genetic material.
S Phase of Interphase Key Events
Cell Growth and Preparation for Cell Division DNA Replication, Increased Protein Synthesis
Similarity to Other Subphases of Interphase All Involve Cell Growth and Preparation for Cell Division
Sister Chromatid Duplication Formation of Identical Chromatids Joined at Centromere
Overall, the S phase of interphase is a critical period of cell growth and preparation for cell division. It involves DNA replication, increased protein synthesis, and the formation of identical sister chromatids that are essential for proper cell division.
G2 Phase of Interphase
The G2 phase is the third and final subphase of interphase, preceded by the G1 phase and S phase. This phase occurs after DNA synthesis and before the onset of cell division through mitosis. The G2 phase serves as a quality control checkpoint to ensure the DNA is intact and ready for cell division.
• The G2 phase is similar to the G1 phase in that the cell is primarily focused on growth and preparing for cell division.
• The DNA is still in the form of chromatin and is not yet condensed into chromosomes.
• The cell undergoes a final check to ensure there are no errors or mutations in the newly synthesized DNA before proceeding with cell division.
Key Event in G2 Phase – DNA Repair and Checkpoint
During the G2 phase of interphase, one of the most important events that occurs is DNA repair and checkpoint. The cell checks the newly synthesized DNA for damage or errors that may prevent proper cell division. The cell has a series of checkpoints to ensure the DNA is intact. If the DNA is damaged, the cell will either repair the damage or undergo apoptosis to prevent the mutation from being passed on to future daughter cells.
The DNA repair process is facilitated by several enzymes, including DNA polymerase, exonucleases, and ligases. These enzymes correct errors in the newly synthesized DNA, ensuring the sequence matches the original DNA template. If the damage is too extensive, the cell may undergo apoptosis to prevent genetic mutations from being propagated.
Checkpoint Protein Function
ATM kinase Monitors DNA damage and activates repair mechanisms
Cyclin-dependent kinase (CDK) Regulates progression through G2 phase and initiates mitosis
BRCA1 and BRCA2 Repair damaged DNA by promoting homologous recombination
The DNA repair process is a complex mechanism that involves many proteins. The checkpoint proteins, ATM kinase, cyclin-dependent kinase (CDK), BRCA1, and BRCA2, play critical roles in ensuring the DNA is intact before cell division proceeds.
DNA Replication during S Phase
The S phase is the second subphase of interphase, and it is characterized by the DNA replication process. This phase is essential as the DNA needs to be replicated before the cell undergoes mitosis. During the S phase, the DNA double helix unwinds, and each strand serves as a template for the synthesis of a new complementary strand of DNA. This process is catalyzed by the enzyme DNA polymerase.
• The S phase occurs after the G1 phase, during which new proteins and organelles are synthesized, and before the G2 phase, which prepares the cell for mitosis.
• The S phase is critical for the accurate transmission of genetic information from parent to daughter cells during cell division.
• DNA replication occurs in a precise and highly controlled manner, ensuring that errors are minimized as much as possible.
The DNA replication process is incredibly complex and involves a multitude of enzymes and proteins. Some of the key events that occur during DNA replication include:
1. Initiation: The replication process is initiated by the formation of a replication bubble, which is mediated by the binding of a protein complex called the origin recognition complex (ORC) to the DNA at specific sites known as origins of replication.
2. Elongation: Once the replication bubble is formed, two replication forks are created, and DNA synthesis proceeds bidirectionally from these forks. Elongation is the process by which DNA polymerase synthesizes new strands of DNA using the parental strands as templates.
3. Termination: The replication process is terminated when the replication forks meet at specific sites known as termination sequences. Once DNA replication is complete, the newly synthesized DNA strands must be checked for errors and repaired before the cell can proceed to mitosis.
DNA Replication during S Phase Description
Initiation The formation of a replication bubble, which is mediated by the binding of a protein complex called the origin recognition complex (ORC) to the DNA at specific sites known as origins of replication.
Elongation The process by which DNA polymerase synthesizes new strands of DNA using the parental strands as templates.
Termination The process by which the replication forks meet at specific sites known as termination sequences, terminating the replication process.
In conclusion, the S phase is a critical subphase of interphase, and DNA replication is the key event that occurs during this phase. DNA replication is a highly precise and controlled process that ensures accurate transmission of genetic information from parent to daughter cells during cell division.
Mitotic Cell Division
Mitotic cell division is the process where a single cell divides into two daughter cells that are genetically identical to each other and the parent cell. This process is crucial in the growth and development of multicellular organisms, as well as in the repair of damaged tissues. Mitotic cell division involves several phases, and one of the key events happens during the S phase of interphase.
The S Phase of Interphase
The S phase is the second subphase of interphase. During this phase, the DNA in the cell is replicated, so that each daughter cell will have a complete copy of the genetic material. This process begins at specific regions of DNA called origins of replication and proceeds in both directions until the entire genome has been duplicated. Once the replication is complete, the cell proceeds to the G2 phase of interphase, where it prepares for mitotic cell division.
Key Events in Mitotic Cell Division
• Prophase: During prophase, the chromatin condenses into visible chromosomes and the nuclear envelope breaks down to allow access to the spindle apparatus.
• Metaphase: In metaphase, the chromosomes align at the equator of the spindle apparatus, called the metaphase plate.
• Anaphase: During anaphase, the sister chromatids are separated and pulled to opposite poles of the cell by the spindle apparatus.
• Telophase: Telophase is the final stage of mitosis, where the nuclear envelope reforms around the two sets of chromosomes, and the cytoplasm divides to produce two daughter cells.
The Mitotic Spindle Apparatus
The mitotic spindle apparatus is a complex network of microtubules that span the cell during mitotic cell division. The spindle apparatus is responsible for segregating the chromosomes into the two daughter cells by attaching to the kinetochores on the centromeres of the chromosomes. The spindle apparatus is composed of three main types of microtubules: astral, polar, and kinetochore microtubules. The spindle apparatus is crucial for the proper distribution of genetic material during mitotic cell division.
Mitotic Phase Key Events
Prophase Condensation of chromatin, breakdown of nuclear envelope, spindle apparatus formation
Metaphase Alignment of chromosomes at the metaphase plate
Anaphase Sister chromatids separate and are pulled to opposite poles of the cell
Telophase Reformation of nuclear envelope, cytoplasm division
The mitotic cell division is a complex process that involves several sub-phases and key events. The S phase of interphase is crucial in the duplication of genetic material. During mitotic cell division, the proper functioning of the mitotic spindle apparatus is necessary for the proper distribution of chromosomes. Understanding the details of mitotic cell division is essential for fields such as medicine and agriculture, as it can help in the understanding and treatment of diseases and the development of new crops.
Regulation of Cell Cycle
The process of cell cycle regulation involves a complex network of pathways and mechanisms that ensure the timely and accurate progression of cells through the various stages of the cell cycle. Dysregulation of these pathways can lead to various diseases, including cancer.
Similarities Among the Three Subphases of Interphase
• All three subphases occur before cell division.
• All three subphases involve DNA replication or synthesis.
• All three subphases play crucial roles in the regulation of the cell cycle.
Key Event During the S Phase
The S phase is the second subphase of interphase and is characterized by DNA replication. During this phase, the genetic material of the cell is duplicated, ensuring that each daughter cell receives an identical set of chromosomes. The key event of the S phase is the replication of the DNA molecule, a process that involves the separation of the double-stranded DNA into two single strands, which then serve as templates for the synthesis of new complementary strands.
Regulation of the Cell Cycle
The regulation of the cell cycle is a tightly controlled process that involves several checkpoints, where the cell must pass specific barriers before proceeding to the next phase. These checkpoints ensure that the DNA is properly replicated, the cell is large enough to undergo division, and the genetic material of the cell is intact and undamaged. If any of these checkpoints are not passed, the cell cycle will be arrested, and the cell will either undergo repair mechanisms or undergo programmed cell death, known as apoptosis.
Checkpoint Function
G1 Checkpoint Ensures that the cell is sufficiently large and has adequate resources to proceed to the S phase.
G2 Checkpoint Verifies the success of DNA synthesis and ensures that the cell is prepared for mitosis.
Mitotic Checkpoint Ensures that the chromosomes are properly attached to the spindle fibers before the cell proceeds with chromosome segregation.
The regulation of the cell cycle is essential for the proper growth and development of an organism. Dysregulation can lead to various diseases, including cancer, where the uncontrolled proliferation of cells results in the formation of a tumor. Understanding the mechanisms of cell cycle regulation is, therefore, crucial for the development of new therapies and treatments for various diseases.
Role of Cyclins and Cyclin-Dependent Kinases (CDKs)
During interphase, cells undergo growth and preparation for replication. The three subphases of interphase include G1 phase, S phase, and G2 phase. While these subphases differ in their activities and processes, the overall goal of interphase is to prepare the cell for division. One crucial component of interphase is the role of cyclins and cyclin-dependent kinases (CDKs).
• Cyclins are proteins that regulate the progression of cells through the cell cycle. They are synthesized during the G1 phase and are then rapidly degraded during the S phase and G2 phase.
• CDKs are enzymes that are activated by cyclins and help regulate the cell cycle. CDKs initiate and drive the events that occur during interphase and mitosis.
• The activity of CDKs is dependent on the presence of the appropriate cyclin. Once cyclins bind to CDKs, they activate the kinase activity of the enzymes, triggering the downstream events necessary for progression through the cell cycle.
During the S phase of interphase, DNA replication occurs, during which DNA is duplicated. The duplicated chromosomes consist of sister chromatids, which are held together by proteins called cohesins. Following DNA replication, the cell enters the G2 phase, where the duplicated chromosomes are checked for errors and repaired if needed before progressing to mitosis.
One of the key events that takes place during S phase is the synthesis of new DNA. The synthesis of new DNA is initiated by CDKs, which in turn rely on the presence of certain cyclins. Specifically, CDK2 and cyclin E play crucial roles in the replication of DNA during S phase.
Role of Cyclins in the Cell Cycle Examples of Cyclins in Humans
Regulate cell cycle progression Cyclin D, Cyclin E, Cyclin A, Cyclin B
Control progression from G1 to S phase Cyclin D
Initiate DNA replication during S phase Cyclin E
Control progression from G2 to M phase Cyclin B
In summary, the role of cyclins and CDKs is critical for proper cell cycle progression during interphase. Cyclins control the timing and order of events that occur during interphase and mitosis, while CDKs initiate and drive the processes necessary for replication and division. During the S phase of interphase, cyclin E and CDK2 are required for the initiation of DNA replication, which is a key event in the preparation for cell division.
Significance of Interphase in Cell Cycle
Interphase is a crucial stage in the cell cycle where the cell prepares for division, and it is made up of three subphases: G1, S, and G2. Although these three subphases have distinct characteristics, they have several similarities.
Similarities among the Three Subphases of Interphase
• All three subphases occur before cell division.
• All three subphases are crucial for the proper functioning of the cell cycle.
• All three subphases involve the replication of DNA.
• All three subphases play a role in ensuring the cell divides correctly.
Key Event During the S Phase
The S phase is the second subphase of interphase and is characterized by the replication of the DNA. During this phase, the DNA helix is separated, and each strand serves as a template for new complementary strands to form. This replication ensures each resulting daughter cell receives a complete copy of the original genetic material.
Importance of Interphase in Cell Cycle
Interphase is a vital stage in the cell cycle as it provides time for the cell to prepare for division, ensure its genetic material is replicated correctly, and check for any errors. Without interphase, cell division would result in incomplete or damaged DNA being distributed to the daughter cells, leading to abnormalities and genetic disorders.
Subphase Characteristics
G1 Cell growth and metabolism; checkpoint for DNA damage
S Replication of DNA; ensures each daughter cell receives a complete copy of the original genetic material
G2 Preparation for mitosis; checkpoint for DNA damage and proper replication
Overall, interphase is a crucial phase in the cell cycle, and its subphases share many similarities while also having distinct characteristics. Understanding the significance of interphase can help us understand how cells divide and play a role in the maintenance of our bodies.
How Are the Three Subphases of Interphase Alike What Key Event Happens During the S Phase
1. What are the three subphases of interphase?
The three subphases of interphase are G1 phase, S phase, and G2 phase.
2. Are the three subphases of interphase alike?
Yes, the three subphases of interphase are alike in that they are all stages of the cell cycle where the cell is not dividing.
3. What happens during the G1 phase?
During the G1 phase, the cell undergoes growth and metabolic activities.
4. What happens during the S phase?
The key event during the S phase is DNA replication, where the cell duplicates its genetic material.
5. What happens during the G2 phase?
During the G2 phase, the cell undergoes final preparations before entering into the phase of cell division.
6. How long does interphase typically last?
Interphase can last anywhere from several hours to several days, depending on the specific type of cell.
7. Why is interphase important for cell division?
Interphase is important for cell division because it is the stage during which the cell grows and replicates its DNA, preparing for the actual process of division.
Closing Thoughts
Thank you for taking the time to read about the key events and similarities in the three subphases of interphase, as well as the importance of the S phase in DNA replication. Remember that interphase is a crucial stage in the cell cycle, setting the stage for successful cell division. Please visit again later for more informative content!
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| 0.995149 |
File: Jamfile
package info (click to toggle)
ftjam 2.3.5-4
• links: PTS
• area: main
• in suites: sarge
• size: 800 kB
• ctags: 1,148
• sloc: ansic: 6,898; yacc: 365; sh: 63; makefile: 58
file content (216 lines) | stat: -rw-r--r-- 4,679 bytes parent folder | download | duplicates (2)
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#
# Jamfile to build Jam (a make(1)-like program)
#
# There are no user-serviceable parts in this file.
#
# Put executables in platform-specific subdirectory.
if $(VMS) { LOCATE_TARGET ?= [.binvms] ; }
else if $(MAC) { LOCATE_TARGET ?= :bin.mac ; }
else { LOCATE_TARGET ?= bin.$(OSFULL[1]:L) ; }
# Leave generated source in current directory; it would be nice to use
# these lines below to build the source into the platform-specific
# directory, but getting scan.c to include the right jambase.h is
# hard: with ""'s, it always gets the bootstrap version; with <>'s,
# it won't find the bootstrap version.
# SEARCH_SOURCE ?= $(LOCATE_TARGET) $(DOT) ;
# LOCATE_SOURCE ?= $(LOCATE_TARGET) ;
#
# We have some different files for UNIX, VMS, and NT.
#
if $(NT) { code = execnt.c filent.c pathunix.c ; }
else if $(OS2)
{
# special case for OS/2. When building Jam with GCC/EMX
# we need to use the "fileunix.c" file
#
# when we build it with other toolsets, we use "fileos2.c"
#
code = execunix.c pathunix.c ;
if $(TOOLSET) = EMX
{
CCFLAGS += -D__OS2__ ;
code += fileunix.c ;
}
else
{
code += fileos2.c ;
}
}
else if $(VMS) { code = execvms.c filevms.c pathvms.c ; }
else if $(MAC) { code = execmac.c filemac.c pathmac.c ; }
else { code = execunix.c fileunix.c pathunix.c ; }
# We have to signal jam.h for these
if $(OS) = NT
{
if $(TOOLSET) = MINGW || $(TOOLSET) = LCC
{
CCFLAGS += -DNT ;
}
else
{
CCFLAGS += /DNT ;
}
}
# Do we know yacc?
if $(YACC) { code += jamgram.y ; }
else { code += jamgram.c ; }
#
# Build the jamgram.y from the jamgram.yy
# yyacc is a slippery script that makes grammars a little
# easier to read/maintain.
#
if $(UNIX) && $(YACC)
{
GenFile jamgram.y jamgramtab.h : ./yyacc jamgram.yy ;
}
#
# How to build the compiled in jambase.
#
Main mkjambase : mkjambase.c ;
#
# The guts of the Jamfile: how to build Jam
#
Main jam : jam.c jambase.c ;
LinkLibraries jam : libjam.a ;
GenFile jambase.c : mkjambase Jambase ;
Library libjam.a :
command.c compile.c $(code) expand.c glob.c
hash.c headers.c hdrmacro.c lists.c make.c make1.c newstr.c
option.c parse.c regexp.c rules.c scan.c search.c subst.c
timestamp.c variable.c ;
if $(BINDIR) { InstallBin $(BINDIR) : jam ; }
#
# Distribution making from here on out.
#
ALLSOURCE =
Build.com Build.mpw Jam.html Jambase Jambase.html Jamfile
Jamfile.html Makefile Porting README RELNOTES command.c command.h
compile.c compile.h execcmd.h execmac.c execunix.c execnt.c execvms.c
expand.c expand.h filemac.c filent.c fileos2.c filesys.h fileunix.c
filevms.c glob.c hash.c hash.h hdrmacro.c hdrmacro.h headers.c
headers.h jam.c jam.h jambase.c jambase.h jamgram.c jamgram.h
jamgram.y jamgram.yy jamgramtab.h lists.c lists.h make.c make.h
make1.c mkjambase.c newstr.c newstr.h option.c option.h parse.c
parse.h patchlevel.h pathmac.c pathunix.c pathvms.c regexp.c regexp.h
rules.c rules.h scan.c scan.h search.c search.h subst.c timestamp.c
timestamp.h variable.c variable.h yyacc
INSTALL
common.mk
builds/win32-visualc.mk
builds/win32-borlandc.mk
builds/win32-gcc.mk
;
rule Binary
{
NotFile package ;
Depends package : $(<) ;
DEPENDS $(<) : $(>) ;
switch $(<)
{
case *-win32.zip : Zip-Exe $(<) : $(>) ;
case *-os2.zip : Zip-Exe $(<) : $(>) ;
case *-linux-libc6.tar : GZip-Exe $(<) : $(>) ;
}
}
rule Package
{
NotFile package ;
Depends package : $(<) ;
DEPENDS $(<) : $(>) ;
switch $(<)
{
case *.tar : { Tar-Gz $(<) : $(>) ; Tar-Bz2 $(<) : $(>) ; }
case *.zip : Zip $(<) : $(>) ;
}
}
VERSION = ftjam-2.3.5 ;
actions Tar-Gz
{
ln -s . $(VERSION)
tar cvhf $(<) $(VERSION)/$(>)
rm $(VERSION)
gzip -9 $(<)
}
actions Tar-Bz2
{
ln -s . $(VERSION)
tar cvhf $(<) $(VERSION)/$(>)
rm $(VERSION)
bzip2 -9 $(<)
}
actions Zip
{
zip -9r $(<) $(>)
}
actions Zip-Exe
{
zip -9j $(<) $(LOCATE_TARGET)\jam.exe
}
actions GZip-Exe
{
ln -s $(LOCATE_TARGET)/jam jam
strip jam
tar chf $(<) jam
rm -f jam
gzip -9 $(<)
}
if $(NT)
{
Binary $(VERSION)-win32.zip : $(ALLSOURCE) ;
Package $(VERSION).zip : $(ALLSOURCE) ;
}
else if $(OS2)
{
Binary $(VERSION)-os2.zip : $(ALLSOURCE) ;
Package $(VERSION).zip : $(ALLSOURCE) ;
}
else if $(OS) = LINUX
{
# how can we detect the C library version reliably ??
# for now, this should only be used for convenience
# purposes, until we add .rpm and .deb support in..
Binary $(VERSION)-linux-libc6.tar : jam ;
Package $(VERSION).tar : $(ALLSOURCE) ;
Package $(VERSION).zip : $(ALLSOURCE) ;
}
|
__label__pos
| 0.834438 |
Article
Chemical and microbial changes during the natural fermentation of strawberry tree (Arbutus unedo L.) fruits
Journal of Food Biochemistry (Impact Factor: 0.76). 01/2007; 31:715-725. DOI:10.1111/j.1745-4514.2007.00138.x
ABSTRACT Aguardente de medronho is the name given in Portugal to a spirit made
from the fermented fruit of Arbutus unedo (strawberry tree), a plant grown
in the Mediterranean region. In order to gain a better understanding of the
fermentation process, as it is performed in the farms, a natural fermentation
with wild microbiota was carried out during 36 days, and some physicochemical
and microbiological parameters were studied. The microbial parameters
analyzed were total viable, lactic and acetic acids bacteria, and yeast counts.
The physicochemical parameters monitored were sugars, minerals, ethanol,
organic acids and pH. Yeasts were the main responsible for the fermentation
of the fruits, as the lactic and acetic acids bacteria are absent. As the fermentation
progressed, the sugars increased during the first 2 days and gradually
decreased along the fermentation period. Maintaining the good quality of the
product could contribute to the preservation and valorization of traditional
resources that are of great importance to prevent their disappearance.
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ABSTRACT: Persimmon (Diospyros kaki) is a seasonal fruit with important health benefits. In this study, persimmon use in wine and condiment production was investigated using molecular methods to identify the yeast and acetic acid bacteria (AAB) isolated from the alcoholic fermentation and acetification of the fruit. Alcoholic fermentation was allowed to occur either spontaneously, or by inoculation with a commercial Saccharomyces cerevisiae wine strain, while acetification was always spontaneous; all these processes were performed in triplicates. Non-Saccharomyces yeast species were particularly abundant during the initial and mid-alcoholic fermentation stages, but S. cerevisiae became dominant toward the end of these processes. During spontaneous fermentation, S. cerevisiae Sc1 was the predominant strain isolated throughout, while the commercial strain of S. cerevisiae was the most common strain isolated from the inoculated fermentations. The main non-Saccharomyces strains isolated included Pichia guilliermondii, Hanseniaspora uvarum, Zygosaccharomyces florentinus and Cryptococcus sp. A distinct succession of AAB was observed during the acetification process. Acetobacter malorun was abundant during the initial and mid-stages, while Gluconacetobacter saccharivorans was the main species during the final stages of these acetifications. Four additional AAB species, Acetobacter pasteurianus, Acetobacter syzygii, Gluconacetobacter intermedius and Gluconacetobacter europaeus, were also detected. We observed 28 different AAB genotypes, though only 6 of these were present in high numbers (between 25%-60%), resulting in a high biodiversity index.
Food Microbiology 05/2012; 30(1):98-104. · 3.41 Impact Factor
• [show abstract] [hide abstract]
ABSTRACT: The edible fruits of four wild small trees or shrubs (Arbutus unedo, Crataegus monogyna, Prunus spinosa, and Rubus ulmifolius) traditionally consumed in the Iberian Peninsula were studied to evaluate their potential for human nutrition, considering their content in bioactive compounds. Lipophilic phytochemicals, such as fatty acids and tocopherols, as well as some hydrophilic antioxidants, such as vitamin C and organic acids, were analyzed. In addition, the antioxidant activity, measured as lipid peroxidation inhibition (β-carotene/linoleate and TBARS assays), was evaluated in each fruit. As far as we know, this is the first report relating to bioactive compounds in wild fruits with relation to the lipid peroxidation inhibition. Data revealed that these wild edible fruits are good sources of bioactive compounds as organic acids, vitamin C, tocopherols, and polyunsaturated fatty acids. They could be considered as functional foods or potential sources of bioactive compounds with antioxidant synergism effect, to be included as antioxidant food ingredients or in dietary supplements, mainly Rubus ulmifolius, due to its high content in tocopherols. This study provides useful and relevant information that justify tocopherols influence in the prevention of lipid peroxidation, due to the strong correlation observed (r>0.95) between these lipophilic bioactive compounds and the antioxidant activity.
European Journal of Lipid Science and Technology 05/2013; 115:176-185. · 2.27 Impact Factor
• Source
[show abstract] [hide abstract]
ABSTRACT: Arbutus wild berries have been traditionally collected and appreciated in the Mediterranean region, although it is actually considered an underutilized fruit-tree species. Increasing the scarce knowledge about its nutritional composition and natural production may interest a broad range of scholars, such as ethnobotanists, chemists, nutritionists and anthropologists. The present study aims to provide original data on the nutritional value and the biomass production of wild strawberry-tree fruits, studying the variation of these characteristics in fruits harvested in different years, from two different Spanish areas. Macro and micronutrient composition of mature Arbutus unedo fruits have been analyzed, with particular attention to the content of some bioactive compounds (fiber, vitamin C as ascorbic and dehydroascorbic acids, total phenolics, carotenoids, including lycopene) and the organic acids profile. The contribution to recommended dietary allowances (RDAs) of this exotic fruit has also been calculated. Fruit crop volume per tree has been estimated as well in the wild strawberry-tree populations surveyed to provide a general framework for discussing the agronomic potential of the species. A wide variability in the nutrient composition of strawberry-tree fruits was found which shows that the analysis of many different samples from different origins and seasons are required to provide average reliable data about the chemical composition of wild fruits. From the results obtained, strawberry-tree fruits can be considered a very good source of health promoting compounds as vitamin C and dietary fiber (202.6 mg/100 g and 42.6% minimum contribution to RDAs, respectively). They are also rich in total available carbohydrates, sugars, potassium and secondary metabolites, such as phenolic compounds, being poor in lipids and Na. These results, together with its high production may help to reinforce its consumption, as an alternative to the fruits available in the market or a source of bioactive compounds for dietary supplements or functional foods.
Food Research International 06/2011; 44(5):1244–1253. · 3.01 Impact Factor
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# # ChangeLog for src/common/extraGlobalenvs.ma # # Generated by Trac 1.2 # Mar 3, 2021, 1:22:45 AM Wed, 06 Mar 2013 11:09:52 GMT piccolo [2783] * src/ERTL/ERTL.ma (modified) * src/ERTL/ERTL_semantics.ma (modified) * src/ERTL/ERTLtoERTLptrOK.ma (modified) * src/ERTLptr/ERTLptr.ma (modified) * src/ERTLptr/ERTLptr_semantics.ma (modified) * src/LIN/LIN_semantics.ma (modified) * src/LIN/joint_LTL_LIN.ma (modified) * src/LIN/joint_LTL_LIN_semantics.ma (modified) * src/LTL/LTL_semantics.ma (modified) * src/RTL/RTL.ma (modified) * src/common/ExtraIdentifiers.ma (added) * src/common/ExtraMonads.ma (modified) * src/common/GenMem.ma (modified) * src/common/extraGlobalenvs.ma (modified) * src/joint/Joint.ma (modified) * src/joint/Traces.ma (modified) * src/joint/extra_joint_semantics.ma (added) * src/joint/joint_semantics.ma (modified) * src/joint/semanticsUtils.ma (modified) modified joint_closed_internal_function definition (added condition ... Wed, 06 Feb 2013 12:01:34 GMT garnier [2608] * src/Clight/Cexec.ma (modified) * src/Clight/CexecSound.ma (modified) * src/Clight/Csem.ma (modified) * src/Clight/MemProperties.ma (modified) * src/Clight/frontend_misc.ma (modified) * src/Clight/memoryInjections.ma (modified) * src/Clight/switchRemoval.ma (modified) * src/Clight/toCminorCorrectnessExpr.ma (modified) * src/Clight/toCminorOps.ma (modified) * src/Cminor/Cminor_semantics.ma (modified) * src/RTLabs/RTLabs_semantics.ma (modified) * src/RTLabs/RTLabs_traces.ma (modified) * src/common/ByteValues.ma (modified) * src/common/GenMem.ma (modified) * src/common/Globalenvs.ma (modified) * src/common/Pointers.ma (modified) * src/common/extraGlobalenvs.ma (modified) Regions are no more stored in blocks. block_region now tests the id, ... Thu, 24 Jan 2013 18:52:38 GMT piccolo [2590] * src/ERTL/semantics.ma (modified) * src/ERTLptr/ERTLtoERTLptrOK.ma (modified) * src/ERTLptr/semantics.ma (modified) * src/common/ExtraMonads.ma (modified) * src/common/PositiveMap.ma (modified) * src/common/extraGlobalenvs.ma (modified) * src/joint/Traces.ma (modified) * src/joint/semantics.ma (modified) added monad machineary for ERTL to ERTLptr translation ... Sat, 05 Jan 2013 12:41:13 GMT piccolo [2570] * src/ERTLptr/ERTLtoERTLptr.ma (modified) * src/ERTLptr/ERTLtoERTLptrOK.ma (added) * src/common/ExtraMonads.ma (added) * src/common/extraGlobalenvs.ma (modified) * src/joint/linearise.ma (modified) * src/joint/lineariseProof.ma (modified) * src/joint/semantics.ma (modified) ERTLtoERTLptr in place Tue, 20 Nov 2012 13:28:06 GMT tranquil [2478] * src/common/extraGlobalenvs.ma (modified) unified is_internal_function_of_program and is_internal_function Mon, 19 Nov 2012 17:04:24 GMT piccolo [2476] * src/common/extraGlobalenvs.ma (modified) * src/joint/lineariseProof.ma (modified) fixed commutation lemmas in lineariseProof started proof of main ... Mon, 19 Nov 2012 09:57:08 GMT tranquil [2474] * src/common/extraGlobalenvs.ma (modified) * src/joint/linearise.ma (modified) changed form of a statement Fri, 16 Nov 2012 17:59:24 GMT tranquil [2473] * src/common/extraGlobalenvs.ma (added) * src/joint/Traces.ma (modified) * src/joint/linearise.ma (modified) * src/joint/semantics.ma (modified) put some generic stuff we need in the back end in extraGlobalenvs ...
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High frequency heat treatment
In short, high frequency quenching (induction heating quenching) refers to the quenching treatment that only the necessary parts are heated and cooled. English for induction hardening, abbreviated as "IH"
Here is a brief introduction to its principle.
① Using electromagnetic induction
Electromagnetic induction refers to that when coil B (two coils are in non-contact state) is placed near coil a connected to the power supply, magnetic field will be generated after the power supply is connected. At this moment, current will also be induced in coil B (the same is true when the power supply is cut off). This phenomenon of induced current in coil B is called electromagnetic induction.
② Heating the product with induced current
If the coil B is replaced by the processed object product (metal), the current induced by electromagnetic induction (eddy current) will pass through the product.
Once there is current passing through the product, Joule heat will be generated due to the resistance, and the product will heat itself. Through this way of direct heating, effective heating can be realized.
③ Heat only the necessary depth of the necessary area
The shape of coil a can be made according to the shape of the product part to be heated, that is, only the necessary parts can be heated.
The higher the frequency of the power supply, the more concentrated the induced current is on the surface of the product. (called skin effect). The frequency can be selected according to the shape (size) of the product and the desired hardening depth.
Vacuum Pump vacuum pump and vacuum furnaces Grinding Machine, Cnc Lathe, Sawing Machine vacuum furnace
vacuum furnace vacuum pump,vacuum furnaces vacuum pump,liquid ring vacuum pump
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Order statistic
From Wikipedia, the free encyclopedia
Jump to: navigation, search
Probability density functions of the order statistics for a sample of size n = 5 from an exponential distribution with unit scale parameter.
In statistics, the kth order statistic of a statistical sample is equal to its kth-smallest value.[1] Together with rank statistics, order statistics are among the most fundamental tools in non-parametric statistics and inference.
Important special cases of the order statistics are the minimum and maximum value of a sample, and (with some qualifications discussed below) the sample median and other sample quantiles.
When using probability theory to analyze order statistics of random samples from a continuous distribution, the cumulative distribution function is used to reduce the analysis to the case of order statistics of the uniform distribution.
Notation and examples[edit]
For example, suppose that four numbers are observed or recorded, resulting in a sample of size 4. If the sample values are
6, 9, 3, 8,
they will usually be denoted
where the subscript i in indicates simply the order in which the observations were recorded and is usually assumed not to be significant. A case when the order is significant is when the observations are part of a time series.
The order statistics would be denoted
where the subscript (i) enclosed in parentheses indicates the ith order statistic of the sample.
The first order statistic (or smallest order statistic) is always the minimum of the sample, that is,
where, following a common convention, we use upper-case letters to refer to random variables, and lower-case letters (as above) to refer to their actual observed values.
Similarly, for a sample of size n, the nth order statistic (or largest order statistic) is the maximum, that is,
The sample range is the difference between the maximum and minimum. It is clearly a function of the order statistics:
A similar important statistic in exploratory data analysis that is simply related to the order statistics is the sample interquartile range.
The sample median may or may not be an order statistic, since there is a single middle value only when the number n of observations is odd. More precisely, if n = 2m+1 for some m, then the sample median is and so is an order statistic. On the other hand, when n is even, n = 2m and there are two middle values, and , and the sample median is some function of the two (usually the average) and hence not an order statistic. Similar remarks apply to all sample quantiles.
Probabilistic analysis[edit]
Given any random variables X1, X2..., Xn, the order statistics X(1), X(2), ..., X(n) are also random variables, defined by sorting the values (realizations) of X1, ..., Xn in increasing order.
When the random variables X1, X2..., Xn form a sample they are independent and identically distributed. This is the case treated below. In general, the random variables X1, ..., Xn can arise by sampling from more than one population. Then they are independent, but not necessarily identically distributed, and their joint probability distribution is given by the Bapat–Beg theorem.
From now on, we will assume that the random variables under consideration are continuous and, where convenient, we will also assume that they have a probability density function (that is, they are absolutely continuous). The peculiarities of the analysis of distributions assigning mass to points (in particular, discrete distributions) are discussed at the end.
Probability distributions of order statistics[edit]
In this section we show that the order statistics of the uniform distribution on the unit interval have marginal distributions belonging to the Beta distribution family. We also give a simple method to derive the joint distribution of any number of order statistics, and finally translate these results to arbitrary continuous distributions using the cdf.
We assume throughout this section that is a random sample drawn from a continuous distribution with cdf . Denoting we obtain the corresponding random sample from the standard uniform distribution. Note that the order statistics also satisfy .
Order statistics sampled from a uniform distribution[edit]
The probability of the order statistic falling in the interval is equal to[2]
that is, the kth order statistic of the uniform distribution is a Beta random variable.[2][3]
The proof of these statements is as follows. For to be between u and u + du, it is necessary that exactly k − 1 elements of the sample are smaller than u, and that at least one is between u and u + du. The probability that more than one is in this latter interval is already , so we have to calculate the probability that exactly k − 1, 1 and n − k observations fall in the intervals , and respectively. This equals (refer to multinomial distribution for details)
and the result follows.
The mean of this distribution is k / (n + 1).
The joint distribution of the order statistics of the uniform distribution[edit]
Similarly, for i < j, the joint probability density function of the two order statistics U(i) < U(j) can be shown to be
which is (up to terms of higher order than ) the probability that i − 1, 1, j − 1 − i, 1 and n − j sample elements fall in the intervals , , , , respectively.
One reasons in an entirely analogous way to derive the higher-order joint distributions. Perhaps surprisingly, the joint density of the n order statistics turns out to be constant:
One way to understand this is that the unordered sample does have constant density equal to 1, and that there are n! different permutations of the sample corresponding to the same sequence of order statistics. This is related to the fact that 1/n! is the volume of the region .
Order statistics sampled from an Erlang distribution[edit]
The Laplace transform of order statistics sampled from an Erlang distribution via a path counting method.[4]
The joint distribution of the order statistics of an absolutely continuous distribution[edit]
If FX is absolutely continuous, it has a density such that , and we can use the substitutions
and
to derive the following probability density functions (pdfs) for the order statistics of a sample of size n drawn from the distribution of X:
where
where
Application: confidence intervals for quantiles[edit]
An interesting question is how well the order statistics perform as estimators of the quantiles of the underlying distribution.
A small-sample-size example[edit]
The simplest case to consider is how well the sample median estimates the population median.
As an example, consider a random sample of size 6. In that case, the sample median is usually defined as the midpoint of the interval delimited by the 3rd and 4th order statistics. However, we know from the preceding discussion that the probability that this interval actually contains the population median is
Although the sample median is probably among the best distribution-independent point estimates of the population median, what this example illustrates is that it is not a particularly good one in absolute terms. In this particular case, a better confidence interval for the median is the one delimited by the 2nd and 5th order statistics, which contains the population median with probability
With such a small sample size, if one wants at least 95% confidence, one is reduced to saying that the median is between the minimum and the maximum of the 6 observations with probability 31/32 or approximately 97%. Size 6 is, in fact, the smallest sample size such that the interval determined by the minimum and the maximum is at least a 95% confidence interval for the population median.
Large sample sizes[edit]
For the uniform distribution, as n tends to infinity, the pth sample quantile is asymptotically normally distributed, since it is approximated by
For a general distribution F with a continuous non-zero density at F −1(p), a similar asymptotic normality applies:
where f is the density function, and F −1 is the quantile function associated with F. One of the first people to mention and prove this result was Frederick Mosteller in his seminal paper in 1946.[5] Further research lead in the 1960s to the Bahadur representation which provides information about the errorbounds.
An interesting observation can be made in the case where the distribution is symmetric, and the population median equals the population mean. In this case, the sample mean, by the central limit theorem, is also asymptotically normally distributed, but with variance σ2/n instead. This asymptotic analysis suggests that the mean outperforms the median in cases of low kurtosis, and vice versa. For example, the median achieves better confidence intervals for the Laplace distribution, while the mean performs better for X that are normally distributed.
Proof[edit]
It can be shown that
where
with Zi being independent identically distributed exponential random variables with rate 1. Since X/n and Y/n are asymptotically normally distributed by the CLT, our results follow by application of the delta method.
Dealing with discrete variables[edit]
Suppose are i.i.d. random variables from a discrete distribution with cumulative distribution function and probability mass function . To find the probabilities of the order statistics, three values are first needed, namely
The cumulative distribution function of the order statistic can be computed by noting that
Similarly, is given by
Note that the probability mass function of is just the difference of these values, that is to say
Computing order statistics[edit]
The problem of computing the kth smallest (or largest) element of a list is called the selection problem and is solved by a selection algorithm. Although this problem is difficult for very large lists, sophisticated selection algorithms have been created that can solve this problem in time proportional to the number of elements in the list, even if the list is totally unordered. If the data is stored in certain specialized data structures, this time can be brought down to O(log n). In many applications all order statistics are required, in which case a sorting algorithm can be used and the time taken is O(n log n). More sophisticated methods can reduce the time to O(n).
See also[edit]
Examples of order statistics[edit]
References[edit]
1. ^ David, H. A.; Nagaraja, H. N. (2003). "Order Statistics". Wiley Series in Probability and Statistics. doi:10.1002/0471722162. ISBN 9780471722168.
2. ^ a b Gentle, James E. (2009), Computational Statistics, Springer, p. 63, ISBN 9780387981444 .
3. ^ Jones, M. C. (2009), "Kumaraswamy’s distribution: A beta-type distribution with some tractability advantages", Statistical Methodology 6 (1): 70–81, doi:10.1016/j.stamet.2008.04.001, As is well known, the beta distribution is the distribution of the m’th order statistic from a random sample of size n from the uniform distribution (on (0,1)).
4. ^ Hlynka, M.; Brill, P. H.; Horn, W. (2010). "A method for obtaining Laplace transforms of order statistics of Erlang random variables". Statistics & Probability Letters 80: 9. doi:10.1016/j.spl.2009.09.006.
5. ^ Mosteller, Frederick (1946). "On Some Useful "Inefficient" Statistics". Annals of Mathematical Statistics (Institute of Mathematical Statistics) 17 (4): 377–408. doi:10.1214/aoms/1177730881. Retrieved February 26, 2015.
• Sefling, R. J. (1980). Approximation Theorems of Mathematical Statistics. New York: Wiley. ISBN 0-471-02403-1.
External links[edit]
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How do I get rid of cholesterol deposits on my skin?
How do I get rid of cholesterol deposits on my skin?
Laser ablation: During this procedure, a doctor will focus a laser beam onto the cholesterol deposit to burn away the fatty tissue. Excision: Excision is the surgical removal of cholesterol deposits by cutting them out of the skin.
What are cholesterol deposits on face?
Cholesterol deposits, also known as xanthomas, are fatty deposits that develop under the skin. Although they can occur anywhere on the body, they most commonly build up around the eyes, and they vary in size from truly minuscule to about three inches across.
How do I get rid of cholesterol bumps on my face?
Treatment for cholesterol deposits around your eyes
1. Surgical excision using a very small blade is typically the first option to remove one of these growths.
2. Chemical cauterization uses chlorinated acetic acids and can remove the deposits without leaving much scarring.
3. Cryotherapy used repeatedly can destroy xanthelasma.
Do cholesterol bumps go away?
Cholesterol deposits are very likely to reoccur following removal, especially in people with high cholesterol. Normalizing lipid levels will have almost no effect on existing deposits. However, treating dyslipidemia is essential, because it can reduce the risk of heart problems.
How do you get rid of cholesterol deposits naturally?
10 Natural Ways to Lower Your Cholesterol Levels
1. Basics.
2. Eat monounsaturated fats.
3. Use polyunsaturated fats.
4. Avoid trans fats.
5. Eat soluble fiber.
6. Exercise.
7. Keep a healthy weight.
8. Don’t smoke.
How do I get rid of milia?
A dermatologist may be able to remove milia from under your eyes using one of the following procedures:
1. Deroofing. A sterilized needle carefully removes the milia from under your eyes.
2. Cryotherapy. Liquid nitrogen freezes the milia, destroying them.
3. Laser ablation.
Can lemon juice lower your cholesterol?
Lemon Juice and Cholesterol Levels According to the latest research, lemon juice may help lower cholesterol levels and improve cardiovascular health. These benefits are largely due to the high levels of flavonoids and vitamin C found in the juice.
Can I remove milia by myself?
Home Removal of Milia Mild cases of milia usually go away on their own. However, you can try a few home procedures that won’t damage your skin and might hasten their departure. A good home remedy is to steam your face. You can do this in the bathroom as part of a hot shower, or you can purchase a facial steamer.
What are these tiny white bumps on my face?
Fordyce spots (Sebaceous Glands)
• Leukonychia (Median Nail Dystrophy)
• Vitiligo
• Pityriasis Alba
• Sun Damaged Skin – Solar Elastosis (hypopigmentation)
• Pityriasis Versicolor (Tinea Versicolor)
• Seborrhoeic Dermatitis (Dandruff Eczema)
• What causes small pimples on the face?
Hormonal Changes in Causes of Pimples: Fluctuations in hormones are one of the top reasons for pimples.
• Excess Production of Sebum: Another main reason behind pimples is the production of excess sebum.
• Dead Skin: Our Skin constantly sheds old cells to give way to new growth.
• What causes blood pimples on my face?
Blood-filled pimples happen as a result of picking or popping a pimple. The forced trauma to that area of the skin pushes out puss and sometimes blood.
How to get rid of cholesterol deposits around your eyes?
Using castor oil is also another effective remedy on how to remove cholesterol deposits around eyes. It helps shrink the cholesterol spots thanks to the presence of the ricinoleic acid in it. However, this oil acts best on small spots that are new. Take pure, cold-pressed castor oil and soak 1 cotton ball in it.
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Like HowStuffWorks on Facebook!
How Jet Packs Work
Barriers to Jet Pack Development
Barriers to Jet Pack Development
The Bell Aerosystems rocket belt in action in 1963.
The Bell Aerosystems rocket belt in action in 1963.
Francis Miller/Time Life Pictures/ Getty Images
The failure to develop a useful jet pack is primarily due to the physics of making a human being fly. We're just not aerodynamic creatures. Nothing in the shape of a human creates lift when we're moved through the air. That means that a jet pack has to create all the lift with pure thrust. Generating that much thrust uses up a lot of fuel -- quickly, too. The true barrier to useful jet pack development is the weight of fuel. We have jet packs, but the reality is that they only work for about 30 seconds. Adding more fuel to prolong flight time would make the jet pack heavier, requiring even more fuel. See the dilemma?
The fuel weight issue has severely limited the usefulness of jet packs. Thirty seconds of flight time isn't enough to do anything other than look for a good place to land. This brings up the second major problem with jet packs: safety. Strapping a rocket or jet engine to your back is inherently dangerous. Shooting yourself up into the air, knowing you only have 30 seconds to get back down, makes it even more dangerous. Since every extra pound that you carry reduces flight time even more, there's not much room for back-up safety systems.
The final flaw with jet engines is one most people don't consider when they fantasize about soaring smoothly through the sky: noise. If you've ever been near a jet engine or a large rocket when it was running, you know they are incredibly loud. One of the early proposed military uses for jet packs was for reconnaissance; however, as soon as the U.S. Army realized that any soldier scouting with jet pack would be heard by the enemy literally miles away, they knew it would never work. The noise would be a problem even in applications where you're not in danger of being shot at. Imagine even one jet pack equipped construction worker in a crowded city. The deafening noise would cause a lot of problems for other workers and anyone unfortunate enough to live or work nearby.
Of course, in the more than 50 years since the first jet packs were developed, we've made some amazing advancements in technology. Surely, if we applied our best scientific minds to the problem, we could overcome these flaws and create functional, useful jet packs, right? It's entirely possible that we might, but there's no demand. It turns out there's not much use for jet packs beyond the, "Hey, cool, jet pack!" factor. Just about any use you can think of for a jet pack can be accomplished with a far cheaper and more reliable technology. It's also not very efficient to transport a single person by air. If we need to get someone somewhere through the air, we can use an airplane or a helicopter and take several people or even some extra cargo.
What was the driving force behind the creation of the first jet pack? Who wanted it built and why? Take a look at the next page to find out.
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What does HLA-DQ2 5 positive mean?
2021-01-13 by No Comments
What does HLA-DQ2 5 positive mean?
In addition, there are at least three different versions of the HLA-DQ2 gene. One, known as HLA-DQ2. 5, confers the highest risk for celiac disease3; about 13% of Caucasian residents of the U.S. carry this specific gene. However, people with other versions of HLA-DQ2 also are at risk for celiac disease.
What is HLA-DQ2 positive?
Interpreting HLA Test Results The risk for celiac disease in patients who are positive for DQ2 or DQ8 can be determined based on which molecules are encoded. DQ2. 5 is most highly associated with celiac disease, while DQ8 and DQ2. 2 are associated with a lower risk.
What is half DQ2?
Haplotypes with DQ2-associated variants in either HLA-DQA1 or HLA-DQB1 are referred to as half-DQ2.
What is HLA typing for celiac disease?
Celiac disease is strongly associated with the HLA genetic region. Approximately 90% of celiac patients express the HLA-DQ2 molecule. Most of the DQ2 negative patients express the HLA-DQ8 molecule. Gluten peptides presented by these HLA molecules induce an abnormal mucosal immune response and tissue damage. …
Are you born with celiac disease?
Most people who are diagnosed with celiac disease are adults. So someone who is born with the genetic risk for the condition can have no autoimmune reaction to gluten for many years, and then for some reason, they break that tolerance to eating gluten and start developing symptoms.
Is there a celiac gene?
“The Human Leukocyte Antigen (HLA) genes are linked to many autoimmune diseases such as celiac disease. Everyone has a copy from the mother and a copy from the father. These copies can come in different versions called alleles. The HLA gene alleles that predispose a person to celiac disease are called DQ2 and DQ8.
What does HLA-DQ2 do?
HLA-DQ molecules bind and present peptides to antigen-specific T cells. It is now commonly accepted that HLA-DQ2. 5 can bind and present gluten peptides and that these HLA-DQ-peptide complexes induce inflammatory T cell responses, causing disease. The HLA-DQ2.
Is celiac disease dominant?
The HLA-DQ2 celiac disease susceptibility haplotype can be inherited either in an autosomal dominant or autosomal recessive manner, depending on the specific alleles of these genes the parent has. HLA-DQ8 celiac disease susceptibility haplotype is inherited in an autosomal dominant manner.
Are you born with celiac disease or do you develop it?
Most people who are diagnosed with celiac disease are adults. So someone who is born with the genetic risk for the condition can have no autoimmune reaction to gluten for many years, and then for some reason, they break that tolerance to eating gluten and start developing symptoms. Studies have confirmed this.
Can 23andMe tell me if I have celiac?
You can see whether you may have a slightly increased risk of developing celiac disease based on your genetics in the 23andMe Celiac Disease Genetic Health Risk report. * 23andMe looks at two common variants associated with celiac disease.
How does the HLA DQ association test work?
Aids in the diagnosis of celiac disease. The HLA DQ Association test provides genotyping for detection of HLA-DQ2 (DQA1*05:01 or 05:05 and DQB1*02:01 or 02:02) and HLA-DQ8 (DQB1*03:02). Patients with DQ2, half DQ2 and/or DQ8 are predisposed to celiac disease.
What is the role of HLA DQA1 in squamous cell carcinoma?
HLA-DQA1 plays an important role in esophageal squamous cell carcinoma progression and may be a biomarker for esophageal squamous cell carcinoma diagnosis and prognosis, as well as a potential target for the treatment of patients with esophageal squamous cell carcinoma.
Is the HLA-DQA1 gene associated with gastric cancer?
HLA-DQA1 gene copy number polymorphism is associated with gastric cancer susceptibility, and there is a multiplicative gene-environment interaction between this polymorphism and H.pylori infection in the development of gastric cancer.
Where are HLA DQ2 and DQA1 alleles most common?
DQ2 β-chains combine with α-chains, encoded by genetically linked HLA-DQA1 alleles, to form the cis – haplotype isoforms. These isoforms, nicknamed DQ2.2 and DQ2.5, are also encoded by the DQA1 *0201 and DQA1 *0501 genes, respectively. DQ2 is most common in Western Europe, North Africa and East Africa.
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Colorant-containing curable negative-type composition, color filter using the composition, and method of manufacturing the same
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The invention provides a colorant-containing curable negative-type composition containing at least a colorant, a cross-linking agent, and a photosensitive agent and satisfying the following Condition (1) relevant to the mass of the respective components in the case the composition further contains a binder, and satisfying the following Condition (2) relevant to the mass of the respective components in the case the composition does not contain a binder: Condition (1): wherein the mass of the binder>0, 0.45≦the mass of the colorant/[the mass of the colorant+the mass of the cross-linking agent+the mass of the photosensitive agent+the mass of the binder] ≦1.0 and 2≦[the mass of the cross-linking agent/ the mass of the binder]; and Condition (2): wherein the mass of the binder=0, 0.45≦the mass of the colorant/[the mass of the colorant+the mass of the cross-linking agent+the mass of the photosensitive agent] ≦1.0.
Skip to: Description · Claims · Patent History · Patent History
Description
CROSS-REFERENCE TO RELATED APPLICATION
This application claims priority under 35USC 119 from Japanese Patent Application No. 2005-138481, the disclosures of which is incorporated by reference herein.
FIELD OF THE INVENTION
The present invention relates to a colorant-containing curable negative-type composition for a color filter that is suitable for forming a colored image of a color filter for use with a liquid crystal display device, a solid image pickup element (such as a CCD or a CMOS) and the like, a color filter, and a method for manufacturing the same.
DESCRIPTION OF THE RELATED ART
As methods for manufacturing a color filter for use with a liquid crystal display device or a solid image pickup element, a dyeing methods, a printing methods, an electrodeposition methods, and a pigment dispersion methods are known.
Among these, the pigment dispersion method is a method for manufacturing a color filter by a photolithography method including using a colored radiation-sensitive composition in which pigments are dispersed in various photosensitive compositions. The pigment dispersion method has an advantage of stability against light, heat and the like, because of the useage of the pigments. In addition, because the pigment dispersion method conducts patterning by the photolithography method, it reveals gives high positioning accuracy. Therefore, the pigment dispersion method has been widely used as a method which is suitable for manufacturing color filters for a color displays of large-screens and high-precision.
In order to manufacture a color filter by the pigment dispersion method:, a radiation-sensitive composition is coated on a glass substrate by means of a spin coater, a roll coater, or the like, and dried for to formation of a coating film; the coating film is exposed to light through a mask pattern, and developed for to formation of colored pixels; and this cycle of operation -is repeated for each color.
A negative type photosensitive composition, which uses a photopolymerizable monomer and a photopolymerization initiator in an alkali soluble resin, is conventionally known as a specific example of the sensitive compositions used in the pigment dispersion method (see, for example, Japanese Patent Application Laid-Open (JP-A) Nos. 2-199403, 4-76062, 5-273411, 6-184482, and 7-140654).
In recent years, for applications, such as solid image pickup elements, it has been demanded that the color filter be of higher precision of the color filter has been demanded. However, it is difficult to further improve the resolution with the conventional pigment dispersion system. In addition, there are problems, such as irregular color being caused by coarse particles of the pigment. Therefore, the pigment dispersion method has not been suited for applications, such as solid image pickup elements where extremely fine patterns are required.
In order to solve the problems, examples where a solvent or water soluble dye is used have been conventionally known ((see, for example, Japanese Patent Application Laid-Open (JP-A) Nos. 2002-278056, and 2002-14221).
However, the dye-containing curable composition involves has the following problems (1) to (4):
• (1) Generally, a coloring matter has a low solubility in either an alkali water solutions or an organic solvents, thus it is difficult to obtain a liquid curing composition having the desired spectrum.
• (2) The dyes may often provide an interaction with other components in the curable composition, thus it is difficult to adjust the solubility (developability) of the curing part and the non-curing part (developability).
• (3) When the dye has a low molar absorption coefficient (ε), the dye must be added in a large quantity, thus it is inevitable to reduce that the amounts of the other components in the curable composition, such as the polymerizable compound (monomer), the binder, and the photopolymerization initiator, must be reduced. This presents such problems as the lowering of the curability of the composition, the heat resistance after curing, and the developability of the (non-)curing part being lowered.
• (4) The dyes are generally poor inferior in light fastness and heat resistance, as compared to pigments.
In addition, especially for the application for manufacturing of a color filter for a solid image pickup element, it is required that the film thickness be 1.5 μm or less, unlike the in semiconductor manufacturing applications. Recently in the view point of performance of a device, it is strongly required that the film thickness be 1.0 μm or less, for example 0.5 to 0.8 μm. Therefore, the coloring matter must be added to the curable composition in a large quantity, which presents the same problems as mentioned above.
Because of the problems, it has been difficult in practice to meet the requirements for performance in the practical for use that is related to the in extremely fine, and thin-film colored patterns for high-precision color filters. Therefore, it has been wished to development of a dye and curable composition which can eliminate the problems has been desired.
SUMMARY OF THE INVENTION
The present invention has been made in view of the above circumstances and provides following a colorant-containing curable negative-type composition of the invention.
A first aspect of the invention is to provide a colorant-containing curable negative-type composition comprising; at least a colorant, a cross-linking agent, and a photosensitive agent, and satisfying the following Condition (1) relevant to the mass of the respective components in the case the composition further contains a binder, and satisfying the following Condition (2) relevant to the mass of the respective components in the case the composition does not contain a binder:
Condition (1): wherein the mass of the binder>0,
0.45≦the mass of the colorant/[the mass of the colorant+the mass of the cross-linking agent+the mass of the photosensitive agent+the mass of the binder]≦1.0 and 2≦[the mass of the cross-linking agent/the mass of the binder]; and
Condition (2): wherein the mass of the binder=0,
0.45≦the mass of the colorant/[the mass of the colorant+the mass of the cross-linking agent+the mass of the photosensitive agent]≦1.0.
A second aspect of the invention is to provide a color filter obtained by using the colorant-containing curable negative-type composition of the invention.
A third aspect of the invention is to provide a method of manufacturing a color filter comprising applying the colorant-containing curable negative-type composition of the invention to a support, exposing the composition through a mask, and developing the composition for to form patterns.
DETAILED DESCRIPTION OF THE INVENTION
The present invention can provide a colorant-containing curable negative-type composition which has a high sensitivity, a high resolution, a high heat resistance, and a wide development latitude; yet is free from elution of the dye; is excellent in solvent resistance of the pattern; and with is high in productivity. Also provided is as well as a color filter using the same, and the method of manufacturing the same.
Especially, present invention is to provide a colorant-containing curable negative-type composition excellent in the molar absorption coefficient and color value of a dye, light fastness, heat resistance, pattern formability (developability), and prevention of color contamination/color omission in the manufacturing process, a color filter in form of a thin film and method of manufacturing the same.
Hereinafter, a colorant-containing curable negative-type composition, a color filter, and its manufacturing method of the invention will be described more in detail.
Colorant-Containing Curable Negative-Type Composition
The colorant-containing curable negative-type composition of the invention (hereinafter, may be referred to as composition of the invention) is a colorant-containing curable negative-type composition comprising; at least a colorant, a cross-linking agent, and a photosensitive agent, and satisfying the following Condition (1) relevant to the mass of the respective components in the case the composition further contains a binder, and satisfying the following Condition (2) relevant to the mass of the respective components in the case the composition does not contain a binder:
Condition (1): wherein the mass of the binder>0,
0.45≦the mass of the colorant/[the mass of the colorant+the mass of the cross-linking agent+the mass of the photosensitive agent+the mass of the binder]≦1.0 and 2≦[the mass of the cross-linking agent/the mass of the binder]; and
Condition (2): wherein the mass of the binder=0,
0.45≦the mass of the colorant/[the mass of the colorant+the mass of the cross-linking agent+the mass of the photosensitive agent]≦1.0.
That is, the composition of the invention may optionally contain a binder, and in the case the composition of the invention contains a binder, the composition satisfies the Condition (1) and in the case the composition of the invention contains no binder, the composition satisfies the Condition (2).
Since the composition of the invention has the feature described above, it can present excellent properties in solubility in a solvent, heat resistance, light fastness, molar absorption coefficient (color value) and pattern formability of colorant.
Since the composition of the invention can improve the unexposed part developability and the remaining film ratio of the exposed part, it has good pattern formability and gives a color filter in form of a thin film and excellent in the various properties. Further, since the composition of the invention scarcely has a risk of deterioration of the properties in the manufacturing process, its productivity is high.
As one embodiment of the invention, it is preferable to satisfy one of the following Conditions (3) to (8) as well as the Conditions (1) and (2).
Condition (3): the colorant contains at least one compound selected from an azo-based compound, a xanthene-based compound, a phthalocyanine-based compound, a triarylmethane-based compound, an anthraquinone-based compound, a quinophthalone-based compound, and an anthrapyridone-based compound.
Condition (4): the colorant contains at least one acidic dye.
Condition (5): a color value (molar absorption coefficient/molecular weight) of the colorant is 15 or more.
Condition (6): the colorant contains at least one material selected from a copper-phthalocyanine coloring material, a pyridone-azo coloring material, a pyrazolone-azo coloring material, a triarylmethane coloring material, a quinophthalone-based coloring material, and a xanthene-based compound coloring material.
Condition (7): the photosensitive agent is a photopolymerization initiator.
Condition (8): the photosensitive agent is a photo-acid generator.
The composition of the invention contains at least a colorant, a cross-linking agent, and a photosensitive agent and if necessary, a binder and a solvent.
Practically, in the case the composition of the invention is a radical polymerizable negative-type composition, the composition contains the colorant, a polymerizable compound (a monomer) as the the cross-linking agent and a photopolymerization initiator as the photosensitive agent and if necessary a binder (preferably an alkali-soluble binder) and a solvent. Also, a thermal cross-linking agent as the cross-linking agent may be further contained.
In the case the composition of the invention is an acid-condensable negative-type composition, the composition contains the colorant, a compound crosslinked and condensed by an acid catalyst as the cross-linking agent, and a photo-acid generator as the photosensitive agent and if necessary a binder (preferably an alkali-soluble binder) and a solvent. Also, a thermal cross-linking agent as the cross-linking agent may be further contained.
Colorant
The composition of the invention contains a colorant. In terms of liquid-phase stability with the lapse of time, the colorant of the invention may contain at least one compound selected from an azo-based compound, a xanthene-based compound, a phthalocyanine-based compound, a triarylmethane-based compound, an anthraquinone-based compound, a quinophthalone-based compound, and an anthrapyridone-based compound and at least one acidic dye. As the colorant those having a color value (molar absorption coefficient/molecular weight) of 15 or more are preferable and those having a color value of 35 or more are more preferable. Further, as the colorant, at least one selected from a copper-phthalocyanine coloring material, a pyridone-azo coloring material, a pyrazolone-azo coloring material, a triarylmethane coloring material, a quinophthalone-based coloring material, and a xanthene-based compound coloring material is preferable.
An azomethine-based coloring material one other than the above exemplified materials is also preferable to be used. The composition of the invention may contain two or more these compounds.
Hereinafter, the colorant (coloring compound),in the invention will be described in detail.
Examples of the colorant in the invention may include conventionally known direct dyes, acidic dyes, mordant dyes/acidic mordant dyes, basic dyes, vat dyes, sulfide dyes, azoic dyes, dispersion dyes, reaction dyes, fluorescent brighteners, and other dyes, pigment resin colors, pigments, near infrared ray absorptive coloring materials, or the like described in COLOR INDEX (SOCIETY OF DYES AND COLOURISTS) and DYEING NOTE (issued by SHIKISENSHA CO., LTD.).
Examples of these colorants are C.I. Solvent Blue 25, C.I. Solvent Blue 55, C.I. Solvent Blue 67, C.I. Solvent Blue 68, C.I. Solvent Blue 38, C.I. Solvent Yellow 82, C.I. Solvent Yellow 162, C.I. Solvent Orange 56, C.I. Acid Violet 17, C.I. Acid Violet 49, C.I. Direct Blue 86, and C.I. Solvent Violet 8. They may be used alone or two or more of them may be used in form of a mixture.
Acidic Dyes
Here is a description about the acidic dyes. The acidic dyes are not particularly limited, provided that they have an acid group, such as the sulfonic acid group or the carboxylic acid group, but an appropriate one of them is selected in consideration of all the necessary performances, such as a solubility in organic solvent and developing solution, a salt formability, an absorbance, an interaction with other components in the curable composition, a light fastness, and a heat resistance.
Specific examples of the acidic dyes are given below, but not limited to these. The examples include: acid alizarin violet N; acid black 1, 2, 24 or 48; acid blue 1, 7, 9, 15, 18, 23, 25, 27, 29, 40, 45, 62, 70, 74, 80, 83, 86, 87, 90, 92, 103, 108, 112, 113, 120, 129, 138, 147, 158, 171, 182, 192, or 249; acid chrome violet K; acid Fuchsin; acid green 1, 3, 5, 9, 16, 25, 27, or 50; acid orange 6, 7, 8, 10, 12, 50, 51, 52, 56, 63, 74, or 95; acid red 1, 4, 8, 14, 17, 18, 26, 27, 29, 31, 34, 35, 37, 42, 44, 50, 51, 52, 57, 66, 73, 80, 87, 88, 91, 92, 94, 97, 103, 111, 114, 129, 133, 134, 138, 143, 145, 150, 151, 158, 176, 183, 198, 211, 215, 216, 217, 249, 252, 257, 260, 266, or 274; acid violet 6B, 7, 9, 17, or 19; acid yellow 1, 3, 7, 9, 11, 17, 23, 25, 29, 34, 36, 40, 42, 54, 65, 72, 73, 76, 79, 98, 99, 111, 112, 114, 116, 169, or 243; Food Yellow 3; and derivatives of these dyes.
Among these, as the acidic dyes, the dyes, such as:
Acid Black 24; Acid Blue 7, 23, 25, 29, 62, 83, 86, 87, 90, 92, 108, 138, 158, 249,; Acid Green 3, 5, 9, 16, 25, 27, 50; Acid Orange 8, 51, 56, 63, 74; Acid Red 1, 4, 8, 34, 37, 42, 52, 57, 80, 97, 114, 143, 145, 151, 183, 217; Acid Violet 7; Acid Yellow 17, 23, 25, 29, 34, 40, 42, 72, 76, 99, 111, 112, 114, 116, 169, 243,and derivatives of these dyes are preferable.
In addition, azo, xanthene, and phthalocyanine acidic dyes other than those mentioned above are preferable, and the acidic dyes, such as C.I. Solvent Blue 44, 38, C.I. Solvent Orange 45, Rhodamine B, Rhodamine 110, 3-[(5-chloro-2-phenoxyphenyl)hydrazono]-3,4-dihydro-4-oxo-5-[(phenylsulfonyl)amino]-2,7-Naphthalenedisulfonic acid, and the derivatives of these dyes are preferably used.
As the derivatives of the acidic dyes, compounds which are produced by converting the sulfonic acid of the acidic dyes into a sulfonamide or a sulfonate, and the like can be effectively used.
Atomic group that forms salts with acidic dyes
About For the atomic group for forming a salt with the acidic dyes, there is no limitation, provided that the atomic group is a cationic one which forms a salt with the anion of the acidic dyes. Examples of such an atomic group include cations consisting of Li, Na, K, Rb, Cs, Ag, Mg, Ca, Sr, Ba, Zn, Al, Ni, Cu, Co, Fe, or a nitrogen-containing compound, and the like.
Among them are cations of H, Li, Na, K, Rb, Cs, Ag, Mg, Ca, Sr, Ba, Zn, Al, Ni, Cu, Co and Fe and cations of nitrogen-containing compounds preferable, cations of H, Na, K, Rb, Cs, Ag, Mg, Ca, Sr, Ba, Zn, Al, Cu, and Fe and cations of nitrogen-containing compounds more preferable, and cations of H, Na, K, Mg, Ca, Ba, Sr, Zn, Al, Cu, and Fe and cations of nitrogen-containing compounds even more preferable.
The cations of nitrogen-containing compounds may be selected in consideration of all of the solubility in an organic solvent and water, salt formability, absorbance and color value of dyes, and heat resistance and light fastness as a colorant. In the case of the selection only based on the absorbance and color value, the nitrogen-containing compounds are preferable to have a molecular weight as low as possible, specifically preferable to have a molecular weight of 300 or less, further more preferable to have a molecular weight of 280 or less, and even more preferable to have a molecular weight of 250 or less.
To prevent photo-discoloration and improve the heat resistance of the dye, a nitrogen-containing compound generally known as a discoloration preventing agent may be used and from this point of view, compounds having lower oxidation potential (lower ionization potential), tertiary amine compounds, alicyclic amine compounds, aniline type compounds, and hydrazine type compounds are preferable.
Hereinafter, practical examples of the nitrogen-containing compounds composing the “cations of the nitrogen-containing compounds” will be exemplified. However, the invention should not be limited to these compounds. Here, the cations are the following nitrogen-containing compounds in cationic state by protonation with one or more protons.
Molar ratio (L) between the atomic group forming a salt with an acidic dye and the acidic dye
Here is an explanation of the ratio between the number of moles of an atomic group which forms a salt with an acidic dye and the number of moles of the acidic dye (hereinafter the ratio may be referred to as “L”). The L is a value determining the ratio of the moles of the acidic dye molecules to the moles of the atomic group, which is its counter ion, and can be freely selected in accordance with the salt forming conditions of the acidic dye and atomic group. Specifically, the L is a numerical value in a range of 0<L≦10, and is the number of acid functional groups in the acidic dye. It is selected in consideration of all the factors, such as the solubility in organic solvent and developing solution, the salt forming properties, the absorbance, the interaction with other components in the curable composition, the light fastness, and the heat resistance. When the selection is carried out only from the viewpoint of absorbance, it is preferable for the L to take a numerical value of 0<L≦7; it is more preferable for the L to take a numerical value of 0<L≦6; and it is particularly preferable for the L to take a numerical value of 0<L≦5.
Color Value of Colorant
The colorant is more preferable as it has higher color value (molar absorption coefficient s/molecular weight Mw) since the addition amount to the dye to the resist can be saved. The color value (i/Mw) is preferably 15 or more, more preferably 20 or more, furthermore preferably 30 or more, and even more preferably 35 or more.
Use Concentration
Next, the use concentration of the colorant will be described. The concentration of the colorant in the total solid components of the composition of the invention, that is although it differs depending on the types, “the mass of the colorant/[the mass of the colorant+the mass of the cross-linking agent+the mass of the photosensitive agent+the mass of the binder] (including the case no binder is used (the mass of the binder=0))” is in a range from 0.45 or more and less than 1.0, preferably from 0.45 or more and 0.99 or less, more preferably 0.45 or more and 0.95 or less, and even more preferably 0.45 or more and 0.90 or less.
If “the mass of the colorant/[the mass of the colorant+the mass of the cross-linking agent+the mass of the photosensitive agent+the mass of the binder] (including the case no binder is used (the mass of the binder=0))” is less than 0.45, it becomes impossible to keep sufficient color density in the case of forming a thin film and if it is 1.0 or more, it becomes impossible to provide sufficient curability.
Cross-Linking Agent
Next, the cross-linking agent in the invention will be described. In the invention, “cross-linking agent” means those forming bonds and forming cross-linking structures. The cross-linking agent in the invention includes, for example, a polymerizable compound, which will be described later, a thermal cross-linking agent, and a compound to be crosslinked and condensed by an acid catalyst.
Note that “the mass of the cross-linking agent”,in the conditions (1) and (2) means the total amount of the cross-linking agent and for example, in the case a thermal cross-linking agent is used in combination with a polymerizable compound, it means the total of these compounds.
Polymerizable Compound
The polymerizable compound (hereinafter, referred to as “monomer”) will be described. The monomer is used together with a photopolymerization initiator, which will be described later, in the case the composition of the invention is a radical polymerizable negative-type composition.
As the monomer, a compound which has at least one addition-polymerizable ethylene group, has a boiling point of 100° C. or more under normal pressure, and has an ethylenic unsaturated group, is preferable. Examples thereof include: monofunctional acrylates and methacrylates such as polyethylene glycol mono(meth)acrylate, polypropylene glycol mono(meth)acrylate, phenoxyethyl (meth)acrylate; polyethylene glycol di(meth)acrylates; trimethylol ethane tri(meth)acrylates; neopentyl glycol di(meth)acrylates; pentaerythritol tri(meth)acrylates; pentaerythritol tetra(meth)acrylates; dipentaerythritol penta(meth)acrylates; dipentaerythritol hexa(meth)acrylates; hexanediol (meth)acrylates; trimethylol propane tri(acryloyloxypropyl) ether; tri(acryloyloxyethyl)isocyanulate; compounds obtained by adding ethylene oxides, propylene oxides or the like to multifunctional alcohols, such as glycerin or trimethylol ethane, and then (meth)acrylating the resultant of the reaction; urethane acrylates such as those disclosed in Japanese Patent Application Publication (JP-B) Nos. 48-41708 and 50-6034 or Japanese Patent Application Laid-Open (JP-A) No. 51-37193; polyester acrylates such as those disclosed in Japanese Patent Application Laid-Open (JP-A) No. 48-64183, Japanese Patent Application Publication (JP-B) Nos. 49-43191 and 52-30490; and multifunctional acrylates or methacrylates, such as epoxyacrylates, which are reaction products of epoxy resins and (metha)acrylic acids, and mixtures thereof. Further, examples thereof includes those introduced as light curable monomers and oligomers in Journal of the Adhesion Society of Japan, Vol. 20, No. 7, pp. 300 to 308.
As the monomer, the (meth)acrylic ester monomer is preferable, and the quadrifunctional or more (meth)acrylic ester monomer is particularly preferable.
From the viewpoint of curability and the like, the content of the monomer in the composition of the present invention is preferably 0.1 to 90% by mass, more preferably 1.0 to 80% by mass, and particularly preferably 2.0 to 70% by mass with respect to the solid content of the composition.
Compound Crosslinked and Condensed by Acid Catalyst
In the case the composition of the invention is an acid-condensable negative-type composition, as the cross-linking agent is used a compound crosslinked and condensed by an acid catalyst. The compound crosslinked and condensed by an acid catalyst includes, for example, a melamine compound, which will be described later as a thermal cross-linking agent.
The content of the compound crosslinked and condensed by an acid catalyst in the composition in the invention is so determined as to satisfy the above Conditions (1) and (2) and it is preferably 0.1 to 55% by mass, more preferably 1.0 to 55% by mass, and even more preferably 2.0 to 55% by mass in the total solid matter of the composition.
Thermal Cross-Linking Agent
Next, the thermal cross-linking agent will be described. The invention provides a film with good curability by advantageously promoting the curing reaction of the film by using the composition as compared with a conventional composition, and it is also possible for the invention to obtain a highly advantageously cured film by using a thermal cross-linking agent as a supplementary agent.
Herein, “thermal cross-linking agent” means an agent to form a bond by stimulation of heat, and thereby form a crosslinked structure. The thermal cross-linking agent to be used in the invention is not particularly limited if it can cure a film by crosslinking reaction and may include (a) an epoxy resin; (b) a melamine compound, a guanamine compound, a glycol uryl compound or an urea compound substituted with at least one substituent group selected from a methylol group, an alkoxymethyl, and an acyloxymethyl; and (c) a phenol compound, a naphthol compound, or a hydroxyanthracene compound substituted with at least one substituent group selected from a methylol group, an alkoxymethyl, and an acyloxymethyl and especially a polyfunctional epoxy resin is preferable.
As the (a) epoxy resin, any compounds can be used with no particular restrictions, provided that they have an epoxy group and crosslikability. Examples of these compounds include: low-molecular weight compounds containing a divalent glyidyl group, such as bisphenol-A-diglycidylether, ethyleneglycol diglycidylether, butanediol diglycidylether, hexanediol diglycidylether, dihydroxybiphenyl diglycidylether, phthalic acid diglycidylether, or N,N-diglycidyl aniline; low-molecular weight compounds containing a trivalent glyidyl group, represented by trimethylolpropane triglycidylether, trimethylolphenol triglycidylether, TRISP-PA (trade name, manufactured by Honshu Chemical Industry Co., Ltd.)-triglycidylether, or the like; low-molecular weight compounds containing a tetravalent glycidyl group, represented by pentaerythritol tetraglycidylether, tetramethylol bisphenol-A-tetraglycidylether, or the like; polyvalent low-molecular weight compounds containing a polyvalent glycidyl group, such as dipentaerythritol pentaglycidylether, dipentaerythritol hexaglycidylether, or the like; high-molecular weight compounds containing a glycidyl group, represented by polyglycidyl (meth)acrylate, a 1,2-epoxy-4-(2-oxylanyl) cyclohexane addition product of 2,2-bis(hydroxymethyl)-1-butanol, or the like.
With regard to a number of the methylol group, alkoxymethyl group or acyloxymethyl group which substitutes the (b) compound, the melamine compound is substituted with 2 to 6 of these substituents, and each of the glycoluryl compound, the guanamine compound, and the urea compound is substituted with 2 to 4 of these substituents. It is preferable that the melamine compound is substituted with 5 to 6 of these substituents, and each of the glycoluryl compound, the guanamine compound, and the urea compound is substituted with 3 to 4.
These methylol-group containing compounds can be obtained by heating the alkoxymethyl-group containing compounds in the presence of an acidic catalyst, such as hydrochloric acid, sulfuric acid, nitric acid, or methansulfonic acid, in alcohol. The acyloxymethyl-group containing compound can be obtained by mixing a methylol-group containing compound with acylchloride and mixing them by stirring in the presence of a basic catalyst.
Hereinafter, specific examples of the (b) compound having substituents are mentioned.
Examples of the melamine compound include a hexamethylol melamine, hexamethoxymethyl melamine, compounds in which 1 to 5 methylol groups of hexamethylol melamine are methoxymethylized or mixtures thereof, hexamethoxyethyl melamine, hexaacyloxymethyl melamine, compounds in which 1 to 5 methylol groups of hexamethylol melamine are acyloxymethylized or mixtures thereof, and the like.
Examples of the guanamine compound include tetramethylol guanamine, tetramethoxymethyl guanamine, compounds in which 1 to 3 methylol groups of tetramethylol guanamine are methoxymethylized or mixtures thereof, tetramethoxyethyl guanamine, tetraacyloxymethyl guanamine, compounds in which 1 to 3 methylol groups of tetramethylol guanamine are acyloxymethylized or mixtures thereof, and the like.
Examples of the glycoluryl compound include tetramethylol glycoluryl, tetramethoxymethyl glycoluryl, compounds in which 1 to 3 methylol groups of tetramethylol glycoluryl are methoxymethylized or mixtures thereof, compounds in which 1 to 3 methylol groups of tetramethylol glycoluryl are acyloxymethylized or mixtures thereof, and the like.
Examples of the urea compound include tetramethylol urea, tetramethoxymethyl urea, compounds in which 1 to 3 methylol groups of tetramethylol urea are methoxymethylized or mixtures thereof, tetramethoxyethyl urea, and the like. These may be used alone or in combination of two or more thereof.
The phenolic compound, naphtholic compound, and hydroxyanthracene compound which are categorized as the (c) compound, being substituted by at least one substituent selected from the group consisting of the methylol group, alkoxymethyl group, and acyloxymethyl group suppress intermixing with an overcoated photoresist by a thermal crosslinking, and further enhance a film strength, as are the case with the (b) compound.
It is necessary that one molecule of the compound (c) includes at least two groups of those selected from the group consisting of the methylol groups, alkoxymethyl groups and acyloxymethyl groups. From the viewpoints of thermal crosslinkability and storage stability, it is preferable that the compound (c) is a phenolic compound in which all of the second and fourth positions are substituted.
Further, it is preferable that the naphtholic compound and the hydroxyanthracene compound, which becomes a skeleton of the compound (c), is substituted at all the ortho and para positions thereof by —OH groups.
The third or fifth position of the phenolic compound, which provides a skeleton of the compound (c), may be substituted or unsubstituted. With regard to the naphtholic compound, which provides the skeleton of the compound (c), the positions other than the ortho positions of —OH groups may be substituted or unsubstituted.
The methylol-group containing compounds can be obtained by using a compound with which the ortho or para position (the second or fourth position) relative to the phenolic -OH group is a hydrogen atom as a raw material, and causing it to react with formalin in a presence of a basic catalyst, such as sodium hydroxide, potassium hydroxide, ammonia, or tetraalkyl ammonium hydroxide.
In addition, the alkoxymethyl-group containing compounds can be obtained by heating the methyrol-group containing compounds in a presence of an acidic catalyst, such as hydrochloric acid, sulfuric acid, nitric acid, or methansulfonic acid, in alcohol.
The acyloxymethyl-group containing compounds can be obtained by causing the methyrol-group containing compounds to react with acylchloride in a presence of a basic catalyst.
Examples of the skeleton compound include a phenolic compound, naphthol, and a hydroxyanthracen compound in which ortho or para positions relative to a phenolic —OH group thereof is unsubstituted, and specific examples thereof include a phenol, crezol, isomers thereof, 2,3-xylenol, 2,5-xylenol, 3,4-xylenol, 3,5-xylenol, bisphenols such as bisphenol-A, 4,4′-bishydroxybiphenyl, TRIS P-PA (described above), naphthol, dihydroxynaphthalen, 2,7-dihydroxyanthracen and the like.
Specific examples of the (c) compound include a trimethylol phenol, tri(methoxymethyl)phenol, trimethylol phenol, compounds in which 1 or 2 methylol groups of trimethylol phenol are methoxymethylized, trimethylol-3-crezol, tri(methoxymethyl)-3-crezol, compounds in which 1 or 2 methylol groups of trimethylol-3-crezol are methoxymethylized, dimethylol crezols such as 2,6-dimethylol-4-crezol, tetramethylol bisphenol-A, tetramethoxymethyl bisphenol-A, compounds in which 1 to 3 methylol groups of tetramethylol bisphenol-A are methoxymethylized, tetramethylol-4,4′-bishydroxybiphenyl, tetramethoxymethyl-4,4-bishydroxybiphenyl, hexamethylol species of TRISP-PA (described above), hexamethoxymethyl species TRIS P-PA (described above), compounds in which 1 to 5 methylol groups of hexamethylol species of TRIS P-PA (described above) are methoxymethylized, bishydroxymethyl naphthalenediol, and the like.
Examples of the hydroxyanthracen compound include a 1,6-dihydroxymethyl-2,7-dihydroxyanthracen, and the like.
Examples of the acyloxymethyl-group containing compounds include compounds in which a part or all of the methylol groups of the methylol-group containing compounds are acyloxymethylized.
Among these compounds, a trimethylol phenol, bishydroxymethyl-p-crezol, tetramethylol bisphenol-A, hexamethylol species of TRIS P-PA (described above), or phenolic compounds in which the hexamethylol groups of these compounds are substituted by an alkoxymethyl group and both a methylol group and an alkoxymethyl group. These may be used alone or in combination of two or more thereof.
The content of the thermal cross-linking agent (components (a) to (c)) in the composition of the invention differs depending on the materials and is so determined as to satisfy the above Conditions (1) to (2) and it is preferably 0 to 55% by mass, more preferably 0 to 50% by mass, and even more preferably 0 to 45% by mass in the total solid matter.
Photosensitive Agent
Next, the photosensitive agent to be added to the composition of the invention will be described. The photosensitive agent may include a photopolymerization initiator in the case the composition of the invention is radical polymerizable negative-type composition and a photo-acid generator in the case the composition of the invention is an acid-condensable negative-type composition.
The photosensitive agent is not particularly limited if it can crosslink the cross-linking agent, however it is preferable to be selected in terms of the properties, initiating efficiency, absorption wavelength, availability, and cost. Photopolymerization initiator
Next, the photopolymerization initiator to be contained when the composition of the present invention is a negative composition will be described. The photopolymerization initiator is not particularly limited, provided that it allows the monomers having polymerizability to be polymerized, however, it is preferable that the photopolymerization initiator is selected from the viewpoints of characteristics, initiation efficiency, absorption wavelength, availability, cost, and the like.
Examples of the photopolymerization initiator include a trihalomethyltriazine compound, a benzyldimethylketal compound, an α-hydroxyketone compound, an a aminoketone compound, a phosphine oxide compound, a metalocen compound, an oxime compound, a triallyl imidazole dimer, a benzothiazole compound, a benzophenone compound, an acetophenone compound and its derivative, a cyclopentadien-benzene-ferrous complex and its salt, a halomethyloxadiazole compound, and a 3-aryl-substituted coumarin compound, and the like, and the photopolymerization initiator preferably comprises at least one compound selected from the group consisting of the a -aminoketone compound, the phosphine oxide compound, the metalocen compound, the oxime compound, and the triallyl imidazole dimer.
In addition, it is preferable that the photopolymerization initiator is a compound which does not generate an acid by decomposition.
Examples of the active halogen compound, such as the halomethyloxadiazole compound, include 2-halomethyl-5-vinyl-1,3,4-oxadiazole compound and the like as disclosed in Japanese Patent Application Publication (JP-B) No. 57-6069, 2-trichloromethyl-5-styryl-1,3,4-oxadiazole, 2-trichloromethyl-5-(p-cyanostyryl)-1,3,4-oxadiazole, 2-trichloromethyl-5-(p-methoxystyryl)-1,3,4-oxadiazole, and the like.
Examples of the trihalomethyl-s-triazine compound photopolymerization initiator include vinyl-halomethyl-s-triazine compounds as disclosed in Japanese Patent Application Publication (JP-B) No. 59-1281, 2-(naphtho-1-yl)-4,6-bis-halomethyl-s-triazine compounds and 4-(p-aminophenyl)-2,6-di-halomethyl-s-triazine compounds as disclosed in Japanese Patent Application Laid-Open (JP-A) No. 53-133428, and the like.
Other examples thereof include a 2,4-bis(trichloromethyl)-6-p-methoxystyryl-s-triazine, 2,6-bis(trichloromethyl)-4-(3 ,4-methylenedioxyphenyl)-1,3,5-triazine, 2,6-bis(trichloromethyl)-4-(4-methylphenyl)-1,3,5-triazine, 2,4-bis(trichloromethyl)-6-( 1-p-di methylaminophenyl-1,3-butadienyl)-s-triazine, 2-trichloromethyl-4-amino-6-p-methoxystyryl-s-triazine, 2-(naphtho-1-yl)-4,6-bis-trichloromethyl-s-triazine, 2-(4-ethoxy-naphtho-1-yl)-4,6-bis-trichlor omethyl-s-triazine, 2-(4-buthoxy-naphtho-1-yl)4,6-bis-trichloromethyl-s-triazine, 2-[4-(2-met hoxyethyl)-naphtho-1-yl]4,6-bis-trichloromethyl-s-triazine, 2-[4-(2-ethoxyethyl)-naphtho-1-y 1]4,6-bis-trichloromethyl-s-triazine, 2-[4-(2-buthoxyethyl)-naphtho-1-yl]4,6-bis-trichloromethyl-s-triazine, 2-(2-methoxy-naphtho-1-yl)4,6-bis-trichloromethyl-s-triazine, 2-(6-methoxy-5-methyl-naphtho-2-yl)-4,6-bis-trichloromethyl-s-triazine, 2-(6-methoxy-naphtho-2-yl)-4,6-bis-trichloromethyl-s-triazine, 2-(5-methoxy-naphtho-1-yl)-4,6-bis-trichloromethyl-s-triazine, 2-(4,7-dimethoxy-naphtho-1-yl)-4,6-bis-trichloromethyl-s-triazine, 2-(6-ethoxy-naphtho-2-yl)-4,6-bis-trichloromethyl-s-triazine, 2-(4,5-dimethoxy-naphtho-1-yl)-4,6-bis-trichloromethyl-s-triazine, 4-[p-N,N-di(ethoxycarbonylmethyl)aminophenyl]-2,6-di(tri chloromethyl)-s-triazine, 4-[o-methyl-p-N,N-di(ethoxycarbonylmethyl)aminophenyl]-2,6-di(trichloromethyl)-s-triazine, 4-[p-N,N-di(chloroethyl)aminophenyl]-2,6-di(trichloromethyl)-s-tri azine, 4-[o-methyl-p-N,N-di(chloroethyl)aminophenyl]-2,6-di(trichloromethyl)-s-triazine, 4-(p-N-chloroethylaminophenyl)-2,6-di(trichloromethyl)-s-triazine, 4-(p-N-ethoxycarbonylmethylaminophenyl)-2,6-di(trichloromethyl)-s-triazine, 4-[p-N,N-di(phenyl)aminophenyl]-2,6-di(trichloromethyl)-s-triazine, 4-(p-N-chloroethylcarbonylaminophenyl)-2,6-di(trichloromethyl)-s-triazine, 4-[p-N-(p-methoxyphenyl)carbonylaminophenyl]-2,6-di(trichloromethyl)-s-triazine, 4-[m-N,N-di(ethoxycarbonylmethyl)aminophenyl]-2,6-di(trichloromethyl)-s-triazine, 4-[m-bromo-p-N,N-di(ethoxycarbonylmethyl)aminophenyl]-2,6-di(trichloromethyl)-s-triazine, 4-[m-chloro-p-N,N-di(ethoxycarbonylmethyl)aminophenyl]-2,6-di(trichloromethyl)-s-triazine, 4-[m-fluoro-p-N,N-di(ethoxycarbonylmethyl)aminophenyl]-2,6-di(trichloromethyl)-s-triazine, 4-[o-bromo-p-N,N-di(ethoxycarbonylmethyl)aminophenyl]-2,6-di(trichloromethyl)-s-triazine, 4-[o-chloro-p-N,N-di(ethoxycarbonylmethyl)aminophenyl]-2,6-di(trichloromethyl)-s-triazine, 4-[o-fluoro-p-N,N-di(ethoxycarbonylmethyl)aminophenyl]-2,6-di(trichloromethyl)-s-triazine, 4-[o-bromo-p-N,N-di(chloroethyl)aminophenyl]-2,6-di(trichloromethyl)-s-triazine, 4-[o-chloro-p-N,N-di(chloroethyl)aminophenyl]-2,6-di(trichloromethyl)-s-triazine, 4-[o-fluoro-p-N,N-di(chloroethyl)aminophenyl]-2,6-di(trichloromethyl)-s-triazine, 4-[m-bromo-p-N,N-di(chloroethyl)aminophenyl]-2,6-di(trichloromethyl)-s-triazine, 4-[m-chloro-p-N,N-di(chloroethyl)aminophenyl]-2,6-di(trichloromethyl)-s-triazine, 4-[m-fluoro-p-N,N-di(chloroethyl)aminophenyl]-2,6-di(trichloromethyl)-s-triazine, 4-(m-bromo-p-N-ethoxycarbonylmethylaminophenyl)-2,6-di(trichloromethyl)-s-triazine, 4-(m-chloro-p-N-ethoxycarbonylmethylaminophenyl)-2,6-di(trichloromethyl)-s-triazine, 4-(m-fluoro-p-N-ethoxycarbonylmethylaminophenyl)-2,6-di(trichloromethyl)-s-triazine, 4-(o-bromo-p-N-ethoxycarbonylmethylaminophenyl)-2,6-di(trichloromethyl)-s-triazine, 4-(o-chloro-p-N-ethoxycarbonylmethylaninophenyl)-2,6-di(trichloromethyl)-s-triazine, 4-(o-fluoro-p-N-ethoxycarbonylmethylaminophenyl)-2,6-di(trichloromethyl)-s-triazine, 4-(m-bromo-p-N-chloroethylaninophenyl)-2,6-di(trichloromethyl)-s-triazine, 4-(m-chloro-p-N-chloroethylaminophenyl)-2,6-di(trichloromethyl)-s-triazine, 4-(m-fluoro-p-N-chloroethylarninophenyl)-2,6-di(trichloromethyl)-s-triazine, 4-(o-bromo-p-N-chloroethylaminophenyl)-2,6-di(trichloromethyl)-s-triazine, 4-(o-chloro-p-N-chloroethylaminophenyl)-2,6-di(trichloromethyl)-s-triazine, 4-(o-fluoro-p-N-chloroethylaminophenyl)-2,6-di(trichloromethyl)-s-triazine, and the like.
In addition to these, TAZ series manufactured by Midori Kagaku Co., Ltd., including TAZ-107, TAZ-110, TAZ-104, TAZ-109, TAZ-140, TAZ-204, TAZ-113, TAZ-123, and TAZ-104 (all trade names, manufactured by Midori Kagaku Co., Ltd.); T series manufactured by Panchim Ltd., including T-OMS, T-BMP, T-R, and T-B (all trade names, manufactured by Panchim Ltd.); IRGACUREO series manufactured by Ciba Specialty Chemicals Inc., including IRGACUREO 651, IRGACURE® 184, IRGACURE® 500, IRGACURE® 1000, IRGACURE® 149, IRGACURE® 819, and IRGACURE® 261; DAROCUR® series manufactured by Ciba Specialty Chemicals Inc., including DAROCUR® 11734; 4′-bis(diethylamino)-benzophenone, 2-(o-benzoyloxime)-1-[4-(phenylthio)phenyl]-1,2-octanedione, 2-benzyl-2-dimethylamino-4-morpholinobuthylophenone, 2,2-dimethoxy-2-phenylacetophenone, 2-(o-chlorphenyl)-4,5-diphenylimidazolyl dimer, 2-(o-fluorophenyl)-4,5-diphenylimidazolyl dimer, 2-(p-methoxyphenyl)-4,5-diphenylimidazolyl dimer, 2-(p-dimethoxyphenyl)-4,5-diphenylimidazolyl dimer, 2-(2,4-dimethoxyphenyl)-4,5-diphenylimidazolyl dimer, 2-(p-methylmercaptophenyl)-4,5-diphenylimidazolyl dimer, benzoinisopropylether, and the like can be usefully used.
Examples of the a -aminoketone compound include IRGACURE® series manufactured by Ciba Specialty Chemicals Inc. (such as IRGACURE®907 or IRGACURE® 369), 2-methyl-1-phenyl-2-morpholinopropane-1-on, 2-methyl-1-[4-(hexyl)phenyl]-2-morpholinopropane-1-on, 2-ethyl-2-dimethylamino-1-(4-morpholinophenyl)-butanone-1, and the like.
The oxime compound is not particularly limited, however, preferable examples include 2-(O-benzoyloxime)-1-[4-(phenylthio)phenyl)-1,2-octanedione, 1-(4-methylsulfanyl-phenyl)-butane-1,2-butane-2-oxime-O-acetate, 1-(4-methylsulfanyl-phenyl)-butane-1-onoxime-O-acetate, hydroxyimino-(4-methylsulfanyl-phenyl)-acetic acid ethyl ester-O-acetate, hydroxyimino-(4-methylsulfanyl-phenyl)-acetic acid ethyl ester-O-benzoate, and the like.
In addition, regarding other photopolymerization initiators, preferable examples of the benzylmethylketal compound include IRGACURE® 65 1; those of the a -hydroxyketone compound include IRGACURE® 184, IRGACURE® 1173, IRGACURE® 500, IRGACURE® 1000, and IRGACURE® 2959; those of the a -aminoketone compound include IRGACURE® 907 and IRGACURE® 369; those of the phosphine oxide compound (blend) include IRGACURE® 1700, IRGACURE® 149, IRGACURE® 1850, IRGACURE® 819, and IRGACURE® 814; those of the metalocen compound include IRGACURE® 784 and IRGACURE® 261 (all manufactured by Cibla Specialty Chemicals Inc.), from the viewpoints of availability and stability, and the analogs/peripheral compounds for these, and the like are also preferable.
As described above, from the viewpoints of light fastness and heat resistance of the dye, it is preferable to use a compound which does not generate an acid by decomposition. Namely, it is preferable to use at least one compound selected from the group consisting of a benzylmethylketal compound, an α-hydroxyketone compound, an α-aminoketone compound, a phosphine oxide compound, a metalocen compound, an oxime compound, a triallylimidazole dimer, a benzothiazole compound, a benzophenone compound, a acetophenone compound and its derivative, and a cyclopentadiene-benzene-ferrous complex and its salt as the compound which does not generate an acid by decomposition. It is further preferable to use at least one compound selected from the group consisting of an α-aminoketone compound, a phosphine oxide compound, a metalocen compound, an oxime compound, and a triallylimidazole dimer.
These photopolymerization initiators can be used in combination with sensitizers and light stabilizers.
Specific examples of the sensitizers and light stabilizers include: benzoin, benzoinmethylether, 9-fluorenone, 2-chloro-9-fluorenone, 2-methyl-9-fluorenone, 9-anthrone, 2-bromo-9-anthrone, 2-ethyl-9-anthrone, 9,10-anthraquinone, 2-ethyl-9,10-anthraquinone, 2-t-butyl-9,10-anthraquinone, 2,6-dichloro-9,10-anthraquinone, xanthone, 2-methylxanthone, 2-methoxyxanthone, 2-methoxyxanthone, thioxanthone, 2,4-diethylthioxanthone, acridone, 10-butyl-2-chloroacridone, benzyl, dibenzalacetone, p-(dimethylamino)phenylstyrylketone, p-(dimethylamino)phenyl-p-methylstyrylketone, benzophenone, p-(dimethylamino)benzophenone(or Michler's ketone), p-(dimethylamino)benzophenone, benzoanthron, and the like; benzothiazole compounds and the like as disclosed in Japanese Patent Application Publication (JP-B) No. 51-48516; TINUVIN® 1130 and TINUVIN® 400 (both manufactured by Ciba Specialty Chemicals); and the like.
In addition to the photopolymerization initiators, commonly-known photopolymerization initiators can be used in the composition of the present invention.
Specific examples thereof include vicinalpolyketolaldonyl compounds as disclosed in U.S. Pat. No. 2,367,660; α-carbonyl compounds as disclosed in U.S. Pat. Nos. 2,367,661 and 2,367,670; acyloinethers as disclosed in U.S. Pat. No. 2,448,828; aromatic acyloin compounds substituted by α-hydrocarbon groups as disclosed in U.S. Pat. No. 2,722,512; multicore quinone compounds as disclosed in U.S. Pat. Nos. 3,046,127 and 2,951,758; combinations of triallylimidazole dimers and p-aminophenylketones as disclosed in U.S. Pat. No. 3,549,367; benzothiazole compounds and trihalomethyl-s-triazine compounds as disclosed in Japanese Patent Application Publication (JP-B) No. 51-48516, and the like.
The amount of the photopolymerization initiator is so determined as to satisfy the Conditions (1) and (2) and it is preferably 0.01% by mass to 50% by mass, more preferably 1% by mass to 40% by mass, and even more preferably 1% by mass to 30% by mass in the solid matters of the monomers. If the amount of the photopolymerization initiator is within a range from 0.01% by mass to 50% by mass, the molecular weight is prevented from becoming too small and accordingly the film strength is prevented from becoming low.
A thermal polymerization inhibitor is preferably added to the composition of the present invention.
Examples of the thermal polymerization inhibitor include hydroquinone, p-methoxyphenol, di-t-butyl-p-crezol, pirrogalol, t-butylcatechol, benzoquinone, 4,4′-thiobis(3-methyl-6-t-butylphenol), 2,2′-methylenebis(4-methyl-6-t-butylphenol), 2-mercaptobenzoimidazole, and the like. Photo-acid generator
In the case the composition of the invention is an acid-condensable negative-type composition [an acid catalyst-curable type (a system in which an acid is generated from a photosensitive agent (a photo-acid generator) by exposure and crosslinking reaction is promoted using the generated acid as a catalyst)], the photo-acid generator is used as the photosensitive agent in combination with the compound (cross-linking agent) to be crosslinked and condensed by the acid catalyst.
The photo-acid generator is not particularly limited if it is a compound generating an acid by exposure and examples may include various oxime-based compounds such as α-(4-toluenesulfonyloxyimino)phenylacetonitrile, various iodonium compounds, various sulfonium compounds, and various trihalomethyltriazine compounds.
The amount of the photo-acid generator is so determined as to satisfy the Conditions (1) and (2), and it is preferably 0.01 to 50% by mass, more preferably 1.0 to 50% by mass, and even more preferably 2.0 to 50% by mass to the solid matter of the compound crosslinked and condensed by the acid catalyst.
Binder
The binder used in the present invention is not particularly limited, provided that it is alkali soluble, however, it is preferable to select the binder from the viewpoints of heat resistance, developability, availability, and the like.
As the alkali soluble binder, a binder which is a linear high polymeric organic substance, is soluble in organic solvent, and can be processed for developement with a weak alkali aqueous solution is preferable. Examptes of such a linear high polymeric organic substance include polymers which have carboxylic acid in its side chain, such as the methacrylic acid copolymer, the acrylic acid copolymer, the itaconic acid copolymer, the crotonic acid copolymer, the maleic acid copolymer, the partially esterified maleic acid copolymer, and the like as disclosed in, for example, Japanese Patent Application Laid-Open (JP-A) Nos. 59-44615, 54-34327, 58-12577, 54-25957, 59-53836, and 59-71048. Similarly, acidic cellulose derivatives which have carboxylic acid in its side chain are useful. In addition to these, a polymer with a hydroxyl group to which an acid anhydride is added, polyhydroxystyrene resins, polysiloxan resins, poly(2-hydroxyethyl (meth)acrylate), polyvinyl pyrrolidone, polyethylene oxide, polyvinyl alcohol, and the like are also useful.
The alkali soluble binder may be a copolymer of monomers having a hydrophilic property, which examples include alcoxyalkyl (meth)acrylate, hydroxyalkyl (meth)acrylate, glycerol (meth)acrylate, (meth)acrylamide, N-methylol acrylamide, secondary or tertiary alkylacrylamide, dialkylaminoalkyl (meth)acrylate, morpholine (meth)acrylate, N-vinyl pirrolidone, N-vinyl caprolactam, vinyl imidazole, vinyl triazole, methyl (meth)acrylate, ethyl (meth)acrylate, branched or straight-chain propyl (meth)acrylate, branched or straight-chain butyl (meth)acrylate, and phenoxyhydroxypropyl(meth)acrylate, and the like.
As other monomers having a hydrophilic property, monomers and the like including a tetrahydrofurfuryl group, phosphoric acid, phosphate ester, quartemary ammonium salt, ethyleneoxy chain, propyleneoxy chain, sulfonic acid and its salt, morpholinoethyl group, and the like are also useful.
Further, in view of improving a crosslinking efficiency, a polymerizable group may be included in the side chain, and polymers and the like which contain an allyl group, a (meta)acryl group, an allyloxyalkyl group, or the like in the side chain thereof are also useful. Examples of the polymers containing these polymerizable groups are given below, but not limited to these, provided that an alkali soluble group, such as α-COOH group, a —OH group, and an ammonium group, and a carbon-carbon unsaturated bond, are included therein.
For example, a compound which is obtained by reacting a compound having an epoxy ring, which has a reactivity with a —OH group, and a compound having a carbon-carbon unsaturated bond group, such as glycidyl acrylate, with a copolymer which is composed of a monomer having a —OH group, such as 2-hydroxyethylacrylate, a monomer having α-COOH group, such as methacrylic acid, and a monomer which is copolymerizable with the monomer having a —OH group and the monomer having α-COOH group, such as an acryl compound, a vinyl compound or the like, can be used. For the reaction with the —OH group, a compound having an acid anhydride, an isocyanate group, and an acryloyl group can be used in place of the epoxy ring. Further, a reaction product which is obtained by reacting a saturated- or unsaturated-polybasic acid anhydride with a compound obtained by reacting a compound having an epoxy ring with an unsaturated carboxylic acid, such as acrylic acid, as disclosed in Japanese Patent Application Laid-Open (JP-A) No. 6-102669 and 6-1938 can also be used.
Examples of a compound which has both an alkali soluble group, such as α-COOH group, and a carbon-to-carbon unsaturated group include DIANAL NR series (trade name, manufactured by Mitsubishi Rayon Co., Ltd.), —COOH group containing polyurethane acrylic oligomer (trade name: PHOTOMER 6173, manufactured by Diamond Shamlock Co., Ltd.), VISCOTE R-264 and KS RESIST 106 (both trade names, manufactured by Osaka Organic Chemical Industry Ltd.), CYCLOMER P series and PRAXEL CF200 series (both trade names, manufactured by Daicel Company Ltd.), EBECRYL 3800 (trade name, manufactured by Daicel-UCB Company Ltd.), and the like.
Among these various binders, preferable examples of the alkali soluble binder to be used in the present invention from the viewpoint of heat resistance include a polyhydroxystyrene resin, polysiloxane resin, (meth)acryl resin, acrylamide resin, and acryl/acrylamide copolymer resin, and particularly preferable examples (meth)acryl resin, polyhydroxystyrene resin, and polysiloxane resin of the alkali soluble binder to be used in the present invention. In addition, from the viewpoint of control of developability, a (meth)acryl resin, acrylamide resin, and acryl/acrylamide copolymer resin are preferable. As the (meth)acrylic resin, a copolymer consisting of monomers selected from a benzyl (meth)acrylate, (meth)acryl acid, hydroxyethyl (meth)acrylate, (meth)acrylamide and the like, (meth)acrylic resins having a polymerizable side-chain, such as CYCLOMER P series, PRAXEL CF200 series (both trade names, manufactured by Daicel Company Ltd.), EBECRYL 3800 (trade name, manufactured by Daicel-UCB Company Ltd.), DIANAL NR series (trade name, manufactured by Mitsubishi Rayon Co., Ltd.), VISCOTE R264, KS RESIST 106 (both trade names, manufactured by Osaka Organic Chemical Industry Ltd.), or the like are preferable.
In addition, in view of enhancing a strength of a cured film, alcohol-soluble nylon, polyether formed from 2,2-bis-(4-hydroxyphenyl)-propane and epichlorhydrine, and the like are also useful.
In addition, examples of the binder to be used in the present invention include an alkali soluble phenolic resin. The alkali soluble phenolic resin can be preferably used when the composition of the present invention is rendered to be a positive composition. Examples of the alkali soluble phenolic resin include a novolak resin, a vinyl copolymer, and the like.
Examples of the novolak resin include a novolak resin which is obtained by condensing phenols and aldehydes in the presence of an acidic catalyst. Examples of the phenols include a phenol, crezol, ethyl phenol, butyl phenol, xylenol, phenyl phenol, catechol, rezorcinol, pyrogallol, naphthol, bisphenol A, and the like. The phenols can be used alone or in combination of two or more of them. Examples of the aldehydes include a formaldehyde, paraformaldehyde, acetaldehyde, propyonic aldehyde, benzaldehyde, and the like.
Specific examples of the novolak resin include a condensation product of a metacrezol, paracrezol, or a mixture of these and formalin. A molecular weight distribution of the novolak resin may be adjusted by means such as fractionation. In addition, a low-molecular weight component having a phenolic hydroxyl group, such as bisphenol C or bisphenol A, may be mixed with the novolak resin.
From the viewpoint of developability, liquid viscosity, and the like, the binder is preferably a copolymer having a weight average molecular weight (a polystyrene-converted value measured by the GPC method) of 1,000 to 2×105, more preferably of 2,000 to 1×105; and particularly preferably of 3,000 to 5×104.
In the case the composition of the invention contains a binder, the content of the binder in the composition of the invention is so determined as to satisfy the Condition (1) and it is preferably 0 to 55% by mass, more preferably 0 to 50% by mass, and even more preferably 0 to 45% by mass to the total solid matter of the composition of the invention.
Also, in the case the composition of the invention contains a binder, it satisfies 2≦[the mass of cross-linking agent/the mass of binder]. The ratio [the mass of cross-linking agent/the mass of binder] satisfies preferably 2≦[the mass of cross-linking agent/the mass of binder]≦∞, more preferably 2.5≦[the mass of cross-linking agent/the mass of binder] <∞, and even preferably 3≦[the mass of cross-linking agent/the mass of binder]≦∞.
Solvent
The solvent to be used with the present invention is not basically particularly limited, provided that the requirements for solubility and coatability of the composition are met, however, it is preferable that the solvent is selected in consideration of the solubility, coatability, and safety of dyes and binders.
Preferable examples of the solvent to be used in preparing the composition of the present invention includes esters such as ethyl acetate, n-butyl acetate, isobutyl acetate, amyl formate, isoamyl acetate, isobutyl acetate, butyl propionate, isopropyl butyrate, ethyl butyrate, butyl butyrate, alkyl esters, methyl lactate, ethyl lactate, methyl oxyacetate, ethyl oxyacetate, butyl oxyacetate, methoxy methylacetate, methoxy ethylacetate, methoxy butylacetate, ethoxy methylacetate, ethoxy ethylacetate, 3-oxypropionic acid alkyl esters such as methyl 3-oxypropionate or ethyl 3-oxypropionate, methyl 3-methoxypropionate, ethyl 3-methoxypropionate, methyl 3-ethoxypropionate, ethyl 3-ethoxypropionate, methyl 2-oxypropionate, ethyl 2-oxypropionate, propyl 2-oxypropionate, methyl 2-methoxypropionate, ethyl 2-methoxypropionate, propyl 2-methoxypropionate, methyl 2-ethoxypropionate, ethyl 2-ethoxypropionate, methyl 2-oxy-2-methylpropionate, ethyl 2-oxy-2-methylpropionate, methyl 2-methoxy-2-methylpropionate, ethyl 2-ethoxy-2-methylpropionate, methyl pyruvate, ethyl pyruvate, propyl pyruvate, methyl acetoacetate, ethyl acetoacetate, methyl 2-oxobutanate, or ethyl 2-oxobutanate; ethers such as diethylene glycol dimethyl ether, tetrahydrofuran, ethylene glycol monomethyl ether, ethylene glycol monoethyl ether, methyl cellosolve acetate, ethyl cellosolve acetate, diethylene glycol monomethyl ether, diethylene glycol monoethyl ether, diethylene glycol monombutyl ether, propylene glycol methyl ether, propylene glycol methyl ether acetate, propylene glycol ethyl ether acetate, or propylene glycol propyl ether acetate; ketones such as methylethyl ketone, cyclohexane, 2-heptanone, or 3-heptanone; and aromatic hydrocarbons such as toluene and xylene.
Among these, as the solvent to be used in the present invention, methyl 3-ethoxypropionate, ethyl 3-ethoxypropionate, ethyl cellosolve acetate, ethyl lactate, diethylene glycol dimethyl ether, butyl acetate, methyl 3-methoxypropionate, 2-heptanone, cyclohexane, ethyl carbitol acetate, butyl carbitol acetate, propylene glycol methyl ether, propylene glycol methyl ether acetate and the like are more preferable.
Variety of additives, such as fillers, high-molecular compounds other than the ones, surfactants, adherence promotors, oxidization inhibitors, ultraviolet absorbers, aggregation inhibitors and the like can be compounded to the composition of the present invention.
Specific examples of these additives include fillers such as glasses or alumina; high-molecular compounds other than binder resins such as polyvinyl alcohol, polyacrylic acid, polyethylene glycol monoalkyl ether, or polyfluoroalkyl acrylate; surfactants such as nonionic surfactants, cationic surfactants, or anionic surfactants; adhesion promotors such as vinyl trimethoxysilane, vinyl triethoxysilane, vinyl tris(2-methoxyethoxy)silane, N-(2-aminoethyl)-3-arninopropylmethyldimethoxysilane, N-(2-aminoethyl)-3-aminopropyltrimethoxysilane, 3-aminopropyltriethoxysilane, 3-glycidoxypropyltrimethoxysilane, 3-glycidoxypropylmethyldimethoxysilane, 2-(3,4-epoxycyclohexyl)ethyltrimethoxysilane, 3-chloropropylmethyldimethoxysilane, 3-chloropropyltrimethoxysilane, 3-methacryloxypropyltrimethoxysilane, or 3-mercaptopropyltrimethoxysilane; oxidation inhibitors such as 2,2-thiobis(4-methyl-6-t-butylphenol) or 2,6-di-t-butylphenol; ultraviolet absorbers such as 2-(3-t-butyl-5-methyl-2-hydroxyphenol)-5-chlorobenzotriazole or alcoxybenzophenone; and aggregation inhibitors such as sodium polyacrylate.
Further, when a promotion of an alkali solubility of a radiation unirradiated portion is intended for further improvement in a developability of the composition of the present invention, an organic carboxylic acid, preferably a low-molecular weight organic carboxylic acid having a molecular weight of 1000 or less, can be added to the composition of the present invention. Specific examples of the organic carboxylic acid include aliphatic monocarboxylic acids such as formic acid, acetic acid, propionic acid, lactic acid, valeric acid, pivalic acid, caproic acid, diethylacetic acid, enanthic acid, or caprylic acid; aliphatic dicarboxylic acids such as oxalic acid, malonic acid, succinic acid, glutaric acid, adipic acid, pimelic acid, suberic acid, azelaic acid, sebacic acid, brasylic acid, methylmalonic acid, dimethylmalonic acid, methylsuccinic acid, tetramethylsuccinic acid, or citraconic acid; aliphatic tricarboxylic acids such as tricarballylic acid, aconitic acid, or camphoronic acid; aromatic monocarboxylic acids such as benzoic acid, toluic acid, cuminic acid, hemellitic acid, or mesitylenic acid; aromatic polycarboxylic acids such as phthalic acid, isophthalic acid, terephthalic acid, trimellitic acid, trimesic acid, mellophanic acid, or pyromellitic acid; and other carboxylic acids such as phenylacetic acid, hydratropic acid, hydrocinnamic acid, mandelic acid, phenylsuccinic acid, atropic acid, cinnamic acid, methyl cinnamate, benzyl cinnamate, cinnamylidene acetic acid, coumaric acid, or unbellic acid.
The content of the solvent in the composition of the invention is preferably 10 to 90% by mass and more preferably 10 to 80% by mass in the total weight of the composition from a viewpoint of solubility of the components, stability with the lapse of time, and coatability.
Color Filter
The color filter of the invention is manufactured by using the composition of the invention. That is, the color filter of the invention contains at least the colorant, the cross-linking agent, the photosensitive agent, and if necessary the binder.
The thickness of the color filter of the invention is preferably 0.1 μm to 1.5 μm, more preferably 0.2 μm to 1.4 μm, furthermore preferably 0.3 μm to 1.3 μm, and even more preferably 0.4 μm to 1.2 μm from a viewpoint of device fabrication.
Further, the color filter of the invention is preferable to have chromaticity change (ΔE*ab) of 20 or less by heat and light. If the chromaticity change (ΔE*ab) by heat and light is 20 or less, the production suitability and storage stability of a product can be improved. The chromaticity change (ΔE*ab) by heat and light is preferably 0 to 18, more preferably 0 to 15, and even more preferably 0 to 13.
The color filter of the present invention can be manufactured by coating the composition of the present invention on a support by using a coating method, such as a spin coating, a casting coating, or a roll coating, to form a radiation-sensitive composition layer, exposing the resultant to light through a prescribed mask pattern, and developing the resultant with a developing solution to form a colored pattern. Further, the manufacturing method for the color filter of the present invention may include a process for curing the resist pattern by heating and/or exposing to light, in accordance with necessity. Moreover, the process for curing by heating and/or exposing to light may be carried out a plurality of times.
Preferable examples of the radiation used in this case include ultraviolet radiation such as g-ray, h-ray or i-ray.
Examples of the support include soda glass, PYREX® glass, silica glass, and these to which a transparent electrically conductive film is deposited; a photoelectric transducer substrate such as a silicone substrate, which is used in image pickup elements or the like; a complementary metal oxide semiconductor (CMOS), and the like. Black stripes for optically isolating the respective pixels may be formed on these supports.
Further, in view of improving adherence the substrate to the upper layer, preventing substance diffusion, and/or flattening the substrate surface, an undercoating layer may be provided on the support in accordance with necessity.
Any developing solution can be used as the developing solution involved in the manufacturing method for the color filter of the present invention, provided that it is a composition which dissolves the composition of the present invention, but do not dissolve radiation-irradiated portions. Specific examples thereof include a combination of various organic solvents and alkaline aqueous solutions.
Examples of the organic solvents include the solvents which are used in preparing the composition of the present invention.
Examples of the alkaline aqueous solution include an alkaline aqueous solution in which an alkaline compound such as sodium hydroxide, potassium hydroxide, sodium carbonate, sodium silicate, sodium metasilicate, ammonia water, ethylamine, diethylamine, dimethylethanolamine, tetramethylammoniumhydroxide, tetraethylammoniumhydroxide, choline, pyrrole, piperidine or 1,8-diazabicyclo-[5.4.0]-7-undecene is dissolved such that a concentration thereof becomes 0.01 to 10% by mass, preferably 0.01 to 1% by mass. When a developing solution consisting of such an alkaline aqueous solution is used, the color filter is generally cleaned with water after a development.
Further, the color filter of the present invention can be used for solid image pickup elements such as a liquid crystal display device or a CCD, and is particularly suited for use with CCDs and CMOSs which have a high resolution of over one million pixels, and the like. The color filter of the present invention can be used as a color filter to be disposed between the light receiving part of the respective pixels constituting the CCD and condensing microlenses, for example.
EXAMPLES
Hereinafter, the present invention will be more specifically described with examples. However, the present invention is not limited by the following examples, in other words, other forms of the invention may be produced within the spirit or scope of the invention as defined in the appended claims. In the following description, the words “part” and “%” are used based on mass unless otherwise noted.
Example 1
1) Preparation of Resist Solution
The following compositions were mixed and dissolved to prepare a resist solution. Compositions for resist solution
Propyleneglycolmonomethyl ether acetate (PGMEA) 19.00 parts Ethyl lactate 36.00 parts Cyclohexanone 0.87 parts Binder (PGMEA solution containing 41% of 30.51 parts acrylmethacrilate/methacrylic acid copolymer (molar ratio = 65:35)) Dipentaerythritolhexacrilate 12.20 parts Polymerization inhibitor (p-methoxyphenol) 0.0075 parts Fluorine surfactant (trade name: MEGAFACE ® F177P, 0.95 parts 0.2% ethyl lactate solution, manufactured by Dainippon Ink And Chemicals, Inc.) 2-(O-benzoyloxime)-1-[4-(phenylthio)phenyl]-1,2- 0.600 parts octanedion (photopolymerization initiator)
2) Preparation of Glass Substrate having Undercoating Layer
A glass substrate (trade name: CORNING 1737, manufactured by Corning Inc.) was ultrasonically washed with 1% NaOH water, which was then followed by water washing and dehydration baking (for 30 min at 200° C.).
Next, the resist solution was applied to the washed glass substrate to provide a film thickness of 2 μm by a spin coater, and heated to be dried for 1 hr at 220° C. for formation of a cured film to obtain a glass substrate having an undercoating layer.
3) Preparation of Colorant-Containing Resist Solution
The following composition was dissolved by mixing to prepare a colorant-containing resist solution.
[Composition of the Colorant-Containing Resist Solution]
Binder solution (Cyclohexanone solution of 20% (benzyl methacrylate/methacrylic acid) copolymer (mole ratio=70:30) (as a polymer, 0.129 part)) 0.645 parts;
Dipentaerythritol hexacrylate (polymerizable compound) 0.516 parts;
The following colorant (I) 0.70 parts;
2-(O-Benzoyloxime)-1-[4-(phenylthio)phenyll-1,2-octanedione (photopolymerization initiator) 0.056 parts; and
Cyclohexanone (solvent) 4.57 parts.
4) Light Exposure and Development Processing of Colorant-Containing Resist (Image Formation)
The dye resist solution obtained in the above paragraph 3) was applied to the undercoating layer of the glass substrate having an undercoating layer obtained in the above paragraph 2) to provide a film thickness of 1.0 μm by using a spin coater, and prebaked for 120 sec at 120° C.
Next, by using an exposing apparatus, the coated film was irradiated at an amount of exposure of 800 mJ/cm2 with a wavelength of 365 nm through a mask having a thickness of 20 μm. After the exposure, the coated film was processed for development under the conditions of 26° C. and 60 sec by using a developing solution (trade name: CD-2000; concentration: 60%, manufactured by FUJIFILM Arch Co., Ltd.) Thereinafter, the coated film was rinsed for 20 sec with running water, which was then followed by spray drying for image formation.
In the present example, the image formation was verified in the normal way by using an optical microscope and SEM photographic observation.
Evaluation
(1) Developability of the Unexposed Parts and Percentage of Film Remaining in the Exposed Parts
The “developability of the unexposed parts” and ” percentage of film remaining in the exposed part” was measured by a color meter (trade name: MCPD-1000, manufactured by Otsuka Electronics Co.,Ltd.).
The “developability of the unexposed parts” refers to the rate of change in absorbance for the film before and after the development, and for a light sensitive negative composition, the greater the value, the better. Further, the “percentage of film remaining in the exposed part” refers to the ratio of maintained light absorbance for the film before and after development, and for a light sensitive negative composition, the greater the value, the better.
That the developability of the unexposed part and the percentage of film remaining in the exposed part are both high values means that pattern formability is good.
(2) Heat Resistance
In the present example, the “heat resistance” was determined by heating the glass substrate coated with a dye resist solution for 1 hr at 200° C. by use of a hot plate, and then measuring the change in chromaticity, i.e., the Δ Eab value, with a chromoscope (trade name: MCPD-1000, manufactured by Otsuka Electronics Co., Ltd.). The smaller the Δ Eab value, the higher the heat resistance.
(3) Light Fastness
The “light fastness” was determined by irradiating the glass substrate coated with a dye resist solution with a Xenon lamp at 200,000 lux for 10 hours (equivalent to 2,000,000 lux hours), and then measuring the change in chromaticity, i.e., the Δ Eab value. The smaller the Δ Eab value, the higher the light fastness.
(4) Molar Absorbance Coefficient (ε)
The molar absorbance coefficient (ε) was calculated from the absorbance in methanol. As the color value, the value obtained by dividing the molar absorbance coefficient (ε) by the Mw of the dye, i.e., ε/Mw was used.
Table 1 gives the results of these.
Examples 2 to 8
The image was formed in the same way as in Example 1, and the evaluation was carried out in the same manner except that, the colorant was substituted by the compounds as shown in Table 1 below. The results are given in Table 1 below.
Example 9
An image was formed in the same manner as Example 1, except that the photopolymerization initiator was changed to “TAZ-107 (manufactured by Midori Kagaku Co., Ltd.) in the preparation of the resist solution of 1) in Example 1 and the image was evaluated similarly. The results are shown in Table 1.
Example 10
An image was formed in the same manner as Example 1, except that the photopolymerization initiator was changed to 2-benzyl-2-dimethylamino-4-morpholinobutylophenone in the preparation of the resist solution of 1) in Example 1 and the image was evaluated similarly. The results are shown in Table 1.
Examples 11 to 20
Except that the glass substrate in Examples 1 to 10 was substituted by a silicone wafer substrate, the image was formed in the same way as in Example 1. For the developability of the unexposed part and the film remaining percentage for the exposed part, the same results as in Examples 1 to 10 were obtained.
Examples 11 to 20 use a silicone wafer substrate, and thus they are different from Examples 1 to 10 in the substrate used, but since the colorant resist solution is coated on the undercoating layer throughout Examples 1 to 20, thus there arises no substantial difference, resulting in the same performances having been obtained.
Comparative Example 1
An image was formed in the same manner as Example 1, except that the composition of the colorant-containing resist solution was changed as follows in the preparation of the colorant-containing resist solution in 3) of Example 1 and the image was evaluated similarly. The results are shown in Table 1.
[Composition of the Colorant-Containing Resist Solution of Comparative Example 1]
Binder solution (Cyclohexanone solution of 20% (benzyl methacrylate/methacrylic acid) copolymer (mole ratio=70:30) (as a polymer, 0.258 parts)) 1.29 parts;
Dipentaerythritol hexacrylate (polymerizable compound) 0.387 parts;
The colorant (I) below 0.70 parts;
2-(O-Benzoyloxime)-1-[4-(phenylthio)phenyl]-1,2-octanedione (photopolymerization initiator) 0.056 parts; and
Cyclohexanone (solvent) 4.57 parts.
Example 21
An image was formed in the same manner as Example 1, except that the composition of the colorant-containing resist solution was changed as follows in the preparation of the colorant-containing resist solution in 3) of Example 1 and the image was evaluated similarly. The results are shown in Table 1.
[Composition of the Colorant-Containing Resist Solution]
The colorant (I) below 16 parts;
α-(4-Toluenesulfonyloxyimino)phenylacetonitrile (photo-acid generator) 0.5 parts;
Hexamethoxymethylmelamine (cross-linking agent) 12.5 parts;
Ethyl lactate (solvent) 70 parts; and
p-Cresol novolak resin (molecular weight 5000) (binder) 6.0 parts.
TABLE 1 Developability Percentage of Heat Light of the film remaining resistance fastness unexposed in the exposed ΔE * ab ΔE * ab ε color value content content Colorant parts (%) parts (%) (200° C./1 h) (2 Mlux/1 h) (l · mol−1cm−1) (ε/Mw) ratio (1) ratio (2) Ex. 1 colorant (I) 100 99 2.25 2.35 45000 66.4 0.5 4.0 Ex. 2 mixture of C.I. Solvent 100 98 3.35 3.25 0.5 4.0 Orange 56 and the colorant (I) at 1.0:1.0 (by mass) Ex. 3 mixture of C.I. Solvent 100 97 3.55 3.65 0.5 4.0 Blue 38 and the colorant (I) at 1.0:1.0 (by mass) Ex. 4 mixture of C.I. Solvent 100 95 4.55 5.25 0.5 4.0 Violet 8 and the colorant (I) at 1.0:1.0 (by mass) Ex. 5 mixture of C.I. Solvent 100 93 5.15 5.45 0.5 4.0 Blue 68 and the colorant (I) at 1.0:1.0 (by mass) Ex. 6 mixture of Acid Red 143 100 92 5.35 5.65 0.5 4.0 and the colorant (I) at 1.0:1.0 (by mass) Ex. 7 mixture of C.I. Solvent 100 91 5.45 5.85 0.5 4.0 Yellow 82 and the colorant (I) at 1.0:1.0 (by mass) Ex. 8 mixture of C.I. Solvent 100 95 15.00 3.25 80 0.5 4.0 Yellow 162 and the colorant (I) at 1.0:1.0 (by mass) Ex. 9 colorant (I) 100 99 4.65 4.85 45000 66.4 0.5 4.0 Ex. 10 colorant (I) 100 97 4.15 4.05 45000 66.4 0.5 4.0 Comp. colorant (I) 100 15 6.25 7.15 45000 66.4 0.5 1.5 Ex. 1 Ex. 21 colorant (I) 100 95 5.35 5.25 45000 66.4 0.45 2.1
Note:
The colorants in Table 1 can be classified in the following types, respectively.
Colorant (I): (pyridone)azo-based acidic dye derivative;
C.I. Solvent Orange 56: xanthene-based dye;
C.I. Solvent Blue 38: copper-phthalocyanine-based dye;
C.I. Solvent Violet 8: triarylmethane-based dye;
C.I. Solvent Blue 68: anthraquinone-based dye;
Acid Red 143: anthrapyridone-based dye;
C.I. Solvent Yellow 82: pyrazolone azo-based dye; and
C.I. Solvent Yellow 162: pyridone azo-based dye in the
Furthermore the content ratios (1) and (2) in Table 1 respectively show as follows.
Content ratio (1): the mass of colorant/[the mass of colorant + the mass of cross-linking agent + the mass of the photosensitive agent + the mass of the binder] and
Content ratio (2): the mass of the cross-linking agent/the mass of the binder.
As shown in Table 1, it has conventionally been difficult to simultaneously satisfy the properties shown in Table 1 just like the case of Comparative Examples. However, use of the colorant-containing curable composition according to the invention makes it possible to obtain a curable composition excellent in the heat resistance, light fastness, and color value (ε/Mw) of the colorant, and curability of the composition. Also, at the same time, the developability of the unexposed part and the percentage of film remaining in the exposed parts are improved and thus the composition is found excellent in the pattern formability.
The curable composition of the invention particularly has an improved curability and pattern formability, so that the film thickness can be reduced even in the case a high film absorbance is required.
Further, as is clear from the results of Table 1, the crosslink density in the film is low in the case of Comparative Example 1 in which the content ratio (2) is as low as (the mass of the cross-linking agent/the mass of the binder)=1.5 and the percentage of film remaining in the exposed parts cannot be maintained.
From the findings, it is proved that the colorant-containing curable composition of the invention is superior.
According to the invention, a curable composition excellent in the heat resistance, light fastness, and molar absorption coefficient (color value) of the colorant is obtained by using a colorant-containing curable composition of the invention. Especially, use of a dye with a high color value makes it possible to lower the addition amount of the dye and accordingly improve various properties of the resist. At the same time, the invention provides a colorant-containing curable composition having improved developability of the unexposed part and percentage of film remaining in the exposed parts and accordingly good pattern formability.
The invention also provides a colorant-containing curable composition with high productivity since it is free of a risk of deterioration of various properties in the manufacturing process and accordingly provides thin film type color filter using the composition.
The invention also provides a simple manufacturing method with a high cost performance by using the colorant-containing curable composition.
According to the invention, a colorant-containing curable negative-type composition having the following characteristics is provided and consequently, the composition has accomplished the aim of the invention.
That is, the invention provides:
<1>. A colorant-containing curable negative-type composition comprising; at least a colorant, a cross-linking agent, and a photosensitive agent, and satisfying the following
Condition (1) relevant to the mass of the respective components in the case the composition further contains a binder, and satisfying the following Condition (2) relevant to the mass of the respective components in the case the composition does not contain a binder:
Condition (1): wherein the mass of the binder>0,
0.45≦the mass of the colorant/[the mass of the colorant+the mass of the cross-linking agent+the mass of the photosensitive agent+the mass of the binder]≦1.0 and 2≦[the mass of the cross-linking agent/the mass of the binder]; and
Condition (2): wherein the mass of the binder=0,
0.45≦the mass of the colorant/[the mass of the colorant+the mass of the cross-linking agent+the mass of the photosensitive agent]≦1.0.
<2>. The colorant-containing curable negative-type composition of <1>, wherein the colorant comprises at least one compound selected from an azo-based compound, a xanthene-based compound, a phthalocyanine-based compound, a triarylmethane-based compound, an anthraquinone-based compound, a quinophthalone-based compound, and an anthrapyridone-based compound.
<3>. The colorant-containing curable negative-type composition of <1>, wherein the colorant comprises at least one acidic dye.
<4>. The colorant-containing curable negative-type composition of <1>, wherein a color value (molar absorption coefficient/molecular weight) of the colorant is 15 or more.
<5>. The colorant-containing curable negative-type composition of <1>, wherein a color value (molar absorption coefficient/molecular weight) of the colorant is 35 or more.
<6>. The colorant-containing curable negative-type composition of <1>, wherein the colorant comprises at least one material selected from a copper-.phthalocyanine coloring material, a pyridone-azo coloring material, a pyrazolone-azo coloring material, a triarylmethane coloring material, a quinophthalone-based coloring material, and a xanthene-based compound coloring material.
<7>. The colorant-containing curable negative-type composition of <1>, wherein the photosensitive agent is a photopolymerization initiator.
<8>. The colorant-containing curable negative-type composition of <1>, wherein the photosensitive agent is a photo-acid generator.
<9>. The colorant-containing curable negative-type composition of <1>, wherein the cross-linking agent comprises a (meth)acrylic ester-based monomer.
<10>. The colorant-containing curable negative-type composition of <1>, wherein the cross-linking agent comprises a tetra- or more-functional (meth)acrylic ester-based monomer.
<11>. The colorant-containing curable negative-type composition of <1>, wherein the binder comprises an alkali-soluble (meth)acrylic resin.
<12>. The colorant-containing curable negative-type composition of <1>, wherein the binder comprises an alkali-soluble (meth)acrylic resin having a polymerizable side chain.
<13>. The colorant-containing curable negative-type composition of <1>, wherein the photosensitive agent comprises at least one compound selected from the group consisting of a trihalomethyltriazine-based compound, a benzyldimethyl ketal compound, an α-hydroxyketone compound, an a-aminoketone compound, a phosphine oxide-based compound, a metallocene compound, an oxime-based compound, a triarylimidazole dimer, a benzothiazole-based compound, a benzophenone compound, an acetophenone compound and derivatives thereof, a cyclopentadiene-benzene-iron complex and salts thereof, a halomethyloxadiazole compound, and a 3-aryl-substituted cumarin compound.
<14>. The colorant-containing curable negative-type composition of <1>, wherein the photosensitive agent comprises at least one compound which does not generate an acid by decomposition.
<15>. The colorant-containing curable negative-type composition of <1>, wherein the photosensitive agent comprises at least one compound selected from the group consisting of an α-aminoketone compound, a phosphine oxide-based compound, a metallocene compound, an oxime-based compound, and a triarylimidazole dimer.
<16>. The colorant-containing curable negative-type composition of <1>, wherein the cross-linking agent comprises a thermal cross-linking agent.
<17>. A color filter obtained by using the colorant-containing curable negative-type composition of <1>.
<18>. The color filter of <17>, wherein a film thickness of the color filter is 0.1 μm to 1.5 μm.
<19>. The color filter of <17>, wherein a chromaticity change (ΔE*ab) of the color filter by heat and light is 20 or less.
<20>. A method of manufacturing a color filter comprising applying the colorant-containing curable negative-type composition of <1>to a support, exposing the composition through a mask, and developing the composition for to form patterns.
Claims
1. A colorant-containing curable negative-type composition comprising; at least a colorant, a cross-linking agent, and a photosensitive agent, and satisfying the following Condition (1) relevant to the mass of the respective components in the case the composition further contains a binder, and satisfying the following Condition (2) relevant to the mass of the respective components in the case the composition does not contain a binder:
Condition (1): wherein the mass of the binder>0,
0.45≦the mass of the colorant/[the mass of the colorant+the mass of the cross-linking agent+the mass of the photosensitive agent+the mass of the binder]≦1.0 and 2≦[the mass of the cross-linking agent/the mass of the binder]; and
Condition (2): wherein the mass of the binder=0,
0.45≦the mass of the colorant/[the mass of the colorant+the mass of the cross-linking agent+the mass of the photosensitive agent]≦1.0.
2. The colorant-containing curable negative-type composition of claim 1, wherein the colorant comprises at least one compound selected from an azo-based compound, a xanthene-based compound, a phthalocyanine-based compound, a triarylmethane-based compound, an anthraquinone-based compound, a quinophthalone-based compound, and an anthrapyridone-based compound.
3. The colorant-containing curable negative-type composition of claim 1, wherein the colorant comprises at least one acidic dye.
4. The colorant-containing curable negative-type composition of claim 1, wherein a color value (molar absorption coefficient/molecular weight) of the colorant is 15 or more.
5. The colorant-containing curable negative-type composition of claim 1, wherein a color value (molar absorption coefficient/molecular weight) of the colorant is 35 or more.
6. The colorant-containing curable negative-type composition of claim 1, wherein the colorant comprises at least one material selected from a copper-phthalocyanine coloring material, a pyridone-azo coloring material, a pyrazolone-azo coloring material, a triarylmethane coloring material, a quinophthalone-based coloring material, and a xanthene-based compound coloring material.
7. The colorant-containing curable negative-type composition of claim 1, wherein the photosensitive agent is a photopolymerization initiator.
8. The colorant-containing curable negative-type composition of claim 1, wherein the photosensitive agent is a photo-acid generator.
9. The colorant-containing curable negative-type composition of claim 1, wherein the cross-linking agent comprises a (meth)acrylic ester-based monomer.
10. The colorant-containing curable negative-type composition of claim 1, wherein the cross-linking agent comprises a tetra- or more-functional (meth)acrylic ester-based monomer.
11. The colorant-containing curable negative-type composition of claim 1, wherein the binder comprises an alkali-soluble (meth)acrylic resin.
12. The colorant-containing curable negative-type composition of claim 1, wherein the binder comprises an alkali-soluble (meth)acrylic resin having a polymerizable side chain.
13. The colorant-containing curable negative-type composition of claim 1, wherein the photosensitive agent comprises at least one compound selected from the group consisting of a trihalomethyltriazine-based compound, a benzyldimethyl ketal compound, an α-hydroxyketone compound, an α-aminoketone compound, a phosphine oxide-based compound, a metallocene compound, an oxime-based compound, a triarylimidazole dimer, a benzothiazole-based compound, a benzophenone compound, an acetophenone compound and derivatives thereof, a cyclopentadiene-benzene-iron complex and salts thereof, a halomethyloxadiazole compound, and a 3-aryl-substituted cumarin compound.
14. The colorant-containing curable negative-type composition of claim 1, wherein the photosensitive agent comprises at least one compound which does not generate an acid by decomposition.
15. The colorant-containing curable negative-type composition of claim 1, wherein the photosensitive agent comprises at least one compound selected from the group consisting of an α-aminoketone compound, a phosphine oxide-based compound, a metallocene compound, an oxime-based compound, and a triarylimidazole dimer.
16. The colorant-containing curable negative-type composition of claim 1, wherein the cross-linking agent comprises a thermal cross-linking agent.
17. A color filter obtained by using the colorant-containing curable negative-type composition of claim 1.
18. The color filter of claim 17, wherein a film thickness of the color filter is 0.1 μm to 1.5 μm.
19. The color filter of claim 17, wherein a chromaticity change (ΔE*ab) of the color filter by heat and light is 20 or less.
20. A method of manufacturing a color filter comprising applying the colorant-containing curable negative-type composition of claim 1 to a support, exposing the composition through a mask, and developing the composition for to form patterns.
Patent History
Publication number: 20060257763
Type: Application
Filed: May 11, 2006
Publication Date: Nov 16, 2006
Patent Grant number: 7763401
Applicant:
Inventor: Katsumi Araki (Shizuoka-ken)
Application Number: 11/431,502
Classifications
Current U.S. Class: 430/7.000; 430/270.100
International Classification: G03F 1/00 (20060101);
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Circle or line from dimension + increment ? does it exist already?
I’m often needing to make a circle just a bit bigger than one that’s already there. etc. You know.
Is there a way to grab the dimension then type in the additional increment I want? Like first click gets the dimension, second click adds the circle, and if numbers are entered it adds them to the initially captured dimension?
No.
Rhino is not a parametric modeler driven by changing dimensions.
However, there is a ModifyRadius command that works on Circles and Arcs.
I don’t want it to change parametrically.
I just want tools to grab a dimension on initial creation, grab a dimension then type in the additional increment I want.
Like a two click solution…make a circle, first click grabs the reference dimension, then I can type in an increment, second click makes the circle equal to initial dimension (click one) plus the entered increment. over and done. No parametric changing or history.
Have a look at the ModifyRadius command please.
That is the closest thing Rhino has to your request.
Ill have a look…is there a similar tool for lines or rectangles?
This is like the main thing missing from Rhino for me. It would be great if all tools had it. Like for line, (vertical, four point, reference + increment).
Note inside the ModifyRadius command, at the prompt for the new radius, you can type in ‘radius’ + Enter to launch the nestable Radius command - you can then select any curve and use the radius for the circle.
A macro will simplfy the process:
! _ModifyRadius _Pause _Radius
-Pascal
Is it possible to make the described approach a feature request?
It really would be nice to have in many tools. It would turn many 6-12 click processes in to 2 click processes.
Hello- does the macro not do what you need?
-Pascal
It changes an existing circle, no? I would like to make a new circle.
But also, a reference dimension+ increment approach would be helpful for many tools. I feel as though this is a fundamental concept missing from rhino. I need something like this all the time.
A liitle more coomplex but it should work, for now
‘! _Circle _Pause _Pause _SelLast _ModifyRadius _Pause _Radius’
To match an existing … right now a script would be need to do what you are asking… I’ll have a look.
-Pascal
Thank you, ill try that out.
But what do you think about a non-parametric reference dimension+ increment function? Or even just a tool that captures a dimension and copies it to the text paste buffer? so it can be pasted in any tools dialogue?
@Pibby - see if this does anything you like -
CirclePlus.py (2.6 KB)
To use the Python script use RunPythonScript, or a macro:
_-RunPythonScript "Full path to py file inside double-quotes"
-Pascal
Offset will make a new circle a specified dimension larger or smaller than the original…
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Hey, just saw this…Ill give it a go and let you know.
Thank you!
Am I doing something wrong?
_-RunPythonScript “C:\Users\Philip\Desktop”
You need to include the file name in your path…
Thanks!
Yes, it works well for circles. Its similar to offset but gives the initial dimension.
But something like this for the move tool, square tool and other common tools would be a real benefit in Rhino.
The circle is just an example.
The concept of a reference point + an increment is something I often wish was present in many tools.
Well, I’m not sure what you’re asking for other than that the original “length” is reported in some commands…
Some things to keep in mind when working in Rhino:
Rhino objects are not parametric. They do not have any data stored on them that tells you how they were made nor any dimensions that can be parametrically changed. They also do not have a “center” or “pivot” point assigned to them, tools like Gumball or BoxEdit calculate object “centers” on the fly using the objects’ bounding box.
Rhino objects also do not generally have knowledge of what specific sub-type of object they are - for example there is not a “square” or “rectangle” entity, they are both just closed planar polylines with specific sets of conditions which were established when they were made - i.e. for a square 4 equal sides and 90° corners, etc. Those conditions are then “forgotten” as soon as the object is created, in Properties all you see is that the object is a closed polyline that has four points.
Of course there are some more sharply defined objects such as circles and arcs because their type of geometry requires it, the radius and center properties for example being inherent to their definition.
There are certainly areas where Rhino could give you more info under specific circumstances, for example Extend (curve) could show you the original length of the curve for reference in addition to the last used extension value on the command line; Move could show you the last used distance/direction information with the possibility to re-use them. Gumball already does this to some extent, but unfortunately, only the move distance or scale factor is remembered, not the axis used.
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Power Spectral Density of Sine-phased BOC signals
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FundamentalsFundamentals
Title Power Spectral Density of Sine-phased BOC signals
Author(s) J.A Ávila Rodríguez, University FAF Munich, Germany.
Level Advanced
Year of Publication 2011
The power spectral density of any BOC[math]\displaystyle{ \left(f_s,f_c\right) }[/math] in sine phasing can be expressed using the theory on Multilevel Coded Spreading Symbols (MCS) signals as follows:
PSD BOC Eq 1.png
(B.1) where the superindex e indicates the even case. Moreover, the modulating factor for the even case presents the following form:
PSD BOC Eq 2.png
(B.2) As we can see, the problem to calculate can be reduced into an easier one by means of the following auxiliary function [math]\displaystyle{ \Phi \left(A\right) }[/math]:
PSD BOC Eq 3.png
and also with the auxiliary function , defined as follows:
PSD BOC Eq 4.png
where we have made the change
PSD BOC Eq 5.png
The interesting property about the above defined function is shown in the next relationship:
PSD BOC Eq 6.png
In fact, taking (4)we can see that
PSD BOC Eq 7.png
Thus (3) can be rewritten as follows
PSD BOC Eq 8.png
Combining now (4) and (6) according to (8) we obtain the following expression:
PSD BOC Eq 9.png
And the modulating function simplifies thus to:
PSD BOC Eq 10.png
where we have also taken into account that according to the definition of the BOC modulation in terms of a BCS vector, n is even in the even version. Additionally, since A can also be expressed as A=jB the expression above simplifies to
PSD BOC Eq 11.png
Once a simplified form has been derived for the BCS modulating factor of BOC(fs, fc), we substitute in (1) yielding the well-known expression for the Power Spectral Density (see Binary Offset Carrier (BOC)):
PSD BOC Eq 12.png
Since , the equation (12) can also be expressed as follows:
PSD BOC Eq 13.png
This is the well known expression that we find everywhere in the literature. Now that we have solved the case of the even BOC modulation in sine phasing, we calculate next its odd counterpart. For the case of the odd BOC modulation in sine phasing, we have to derive first a general expression for any odd n. We will proceed by generalizing over n.
For n = 3, [math]\displaystyle{ BOC_{sin}\left(f_s,f_c\right) }[/math] can also be expressed as BCS([+1,-1,+1], fc), such that the generation matrix will adopt the following form:
PSD BOC Eq 14.png
Thus, the odd modulating term yields this time:
PSD BOC Eq 15.png
where o indicates the odd case, and
PSD BOC Eq 16.png
In the same manner, for , [math]\displaystyle{ BOC_{sin}\left(f_s,f_c\right)=BOC_{sin}\left(2f_c,f_c\right) }[/math] what can also be defined as in the general form BCS([+1,-1,+1,-1,+1], fc) with generation matrix given by:
PSD BOC Eq 17.png
Thus, for the case of , we will have:
PSD BOC Eq 18.png
If we continue by induction we can find the expression for any odd n:
PSD BOC Eq 19.png
As we can recognize, (19) is equal to (2) except that n is odd now with [math]\displaystyle{ n \in \left \{ 3,5,7,\cdots \right \} }[/math] Moreover, it can be shown that for the odd [math]\displaystyle{ n=2f_s/f_c }[/math] is still valid. For simplicity, we express the modulating factor above using its exponential equivalent expression:
PSD BOC Eq 20.png
Using now the expressions derived above for the sum term [math]\displaystyle{ \Phi\left(A\right) }[/math], it can be shown that:
PSD BOC Eq 21.png
Again, this expression is similar to that obtained for the even case, but with a slight difference. Indeed, since n is odd, the second summand in the numerator has a changed sign with respect to (9). The modulating function is thus shown to present the following form:
PSD BOC Eq 22.png
Additionally, since A can also be expressed as A=jB, (22) simplifies to
PSD BOC Eq 23.png
Once we have the BCS modulating factor of an arbitrary odd [math]\displaystyle{ BOC\left(f_s,f_c\right) }[/math] it can be shown that the power spectral density is (see Binary Offset Carrier (BOC)):
PSD BOC Eq 24.png
which coincides perfectly with the expressions found in the literature [J. W. Betz, 1999][1], [A.R. Pratt and J.I.R. Owen, 2003][2] and [E. Rebeyrol et al., 2005][3].
Furthermore, since [math]\displaystyle{ n=2f_s/f_c }[/math], the previous expression can also be shown as follows:
PSD BOC Eq 25.png
References
1. ^ [J. W. Betz, 1999] J. W. Betz, The offset carrier modulation for GPS modernization, in Proceedings of the National Technical Meeting of the Institute of Navigation, ION-NTM 1999, pp. 639–648, January 1999, San Diego, California, USA.
2. ^ [A.R. Pratt and J.I.R. Owen, 2003] Anthony R. Pratt & John I.R. Owen, BOC Modulations Waveform, Proceedings of the International Technical Meeting of the Institute of Navigation, ION-GNSS 2003, 9-12 September 2003, Portland, Oregon, USA.
3. ^ [E. Rebeyrol et al., 2005] E. Rebeyrol, C. Macabiau, L. Lestarquit, L. Ries, J-L. Issler, M.L. Boucheret, M. Bousquet , BOC Power Spectrum Densities, Proceedings of the National Technical Meeting of the Institute of Navigation, ION-NTM 2005, 24-26 January 2005, Long Beach, California, USA.
Credits
The information presented in this NAVIPEDIA’s article is an extract of the PhD work performed by Dr. Jose Ángel Ávila Rodríguez in the FAF University of Munich as part of his Doctoral Thesis “On Generalized Signal Waveforms for Satellite Navigation” presented in June 2008, Munich (Germany)
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Does anyone have a elegant way of dealing with errors in ASP.Net MVC? I constantly run into issues when dealing with requests to controller actions where the Action can be used for both normal requests and AJAX requests. The problem I have is finding an elegant way of dealing with these issues.
For example, how could I handle validation errors? Ideally I would like to submit the form to a server via AJAX and then return any errors the action threw and display them on the page, but for the same to work via a normal postback when the client has JavaScript turned off. I know I can use the jQuery Validation Plugin as they type but it's isn't the same, nor is it ideal considering the restrictions on the data to be validated would be specified in two places (my nHibernate Validation Mappings and in the JavaScript file).
How about when the user requests a non-existent record? Should I redirect to a 404 page? What if the request was made via Ajax (say, to load in a dialogue box).
So:
How do you handle errors thrown by controller actions when they were called using Ajax? Especially Model State errors (i.e validation). Can I send it via JSON?
Do you have tips on how to make controller actions which work well when called normally and via Ajax? This is a annoying problem when writing action methods. Because of the return type I may want a different result depending on the caller. Is there anyway to do this without having two Action Methods?
What is your general strategy for handling errors in actions on MVC? Do you redirect to error pages a lot? Do you redirect to another page?
Edit: I think part of the problem is I want different things to happen, so if there is an error I would want to stop any progress until it's fix and therefore send a error back. Otherwise I may want to update different areas of the page. But if I had one return how would I know it's a success or failure without wrapping it an object that has a property indicting so (thus making it more difficult to use partial views)
Thanks
share|improve this question
I'd also like to hear the answer to your second question. I have some major difficulty in the past trying to get this to work out. – Merritt Aug 4 '09 at 17:21
It is possible to detect whether your request is an AJAX call and return appropriate data. See my (rough) example. – Peter J Aug 5 '09 at 14:53
7 Answers 7
up vote 3 down vote accepted
The AJAX call is not a page refresh so I would definitely not redirect to a 403, 404 etc. I would display a client side dialog explaining the unexpected result of the request appropriately. Your controller can return the results of failed validation to the AJAX call in a similar fashion to how it would return success data so your dialog can also display whatever is required in this scenario.
share|improve this answer
In this case the Dialog would close on success but anyway, What would you do if the data returned was different between a success and a error – Damien Aug 4 '09 at 16:26
1
If success, close the dialog, if failure, display an appropriate message? – Adam Ralph Aug 4 '09 at 20:09
1
What AdamRalph suggested is exactly what we do. I built a custom view "_NotificationView" that the modals return, on ajax return I check to see if the html is blank, if it is, I close the dialog, otherwise I display the message to the user and keep the dialog open. – Tom Anderson Aug 5 '09 at 14:39
Do you have tips on how to make controller actions which work well when called normally and via Ajax? This is a annoying problem when writing action methods.
Yes, yes I do. We also worked through a similar problem--we wanted the app to have a bunch of forms which would be called generally via ajax but could get hit normally. Moreover, we didn't want a whole bunch of duplicate logic in javascript floating around. Anyhow, the technique we came up with was to use a pair of ActionFilterAttributes to intercept the forms, then a little javascript to wire up the form for ajax handling.
First, for the ajax request, we just wanted to swap out master pages:
private readonly string masterToReplace;
/// <summary>
/// Initializes an Ajax Master Page Switcharoo
/// </summary>
/// <param name="ajaxMaster">Master page for ajax requests</param>
public AjaxMasterPageInjectorAttribute(string ajaxMaster)
{
this.masterToReplace = ajaxMaster;
}
public override void OnResultExecuting(ResultExecutingContext filterContext)
{
if (!filterContext.HttpContext.Request.IsAjaxRequest() || !(filterContext.Result is ViewResult)) return;
ViewResult vr = (ViewResult) filterContext.Result;
vr.MasterName = masterToReplace;
}
}
On the return side, we use xVal to validate client-side, so one shouldn't get much in the way of invalid data, but one can still get some. To that end, we just use normal validation and have the aciton method return the form with validation messages. Successful posts are rewarded with redirects in general. In any case, we do a little json injection for the success case:
/// <summary>
/// Intercepts the response and stuffs in Json commands if the request is ajax and the request returns a RedirectToRoute result.
/// </summary>
public class JsonUpdateInterceptorAttribute : ActionFilterAttribute
{
public override void OnActionExecuted(ActionExecutedContext filterContext)
{
if (filterContext.HttpContext.Request.IsAjaxRequest())
{
JsonResult jr = new JsonResult();
if (filterContext.Result is RedirectResult)
{
RedirectResult rr = (RedirectResult) filterContext.Result;
jr.Data = new {command = "redirect", content = rr.Url};
}
if (filterContext.Result is RedirectToRouteResult)
{
RedirectToRouteResult rrr = (RedirectToRouteResult) filterContext.Result;
VirtualPathData vpd = RouteTable.Routes.GetVirtualPath(filterContext.RequestContext, rrr.RouteValues);
jr.Data = new {command = "redirect", content = vpd.VirtualPath};
}
if (jr.Data != null)
{
filterContext.Result = jr;
}
}
}
}
The final trick is using a little javascript object to tie everything together:
function AjaxFormSwitcher(form, outputTarget, doValidation) {
this.doValidation = doValidation;
this.outputTarget = outputTarget;
this.targetForm = form;
}
AjaxFormSwitcher.prototype.switchFormToAjax = function() {
var afs = this;
var opts = {
beforeSubmit: this.doValidation ? afs.checkValidation : null,
complete: function(xmlHttp, status){ afs.processResult(afs, xmlHttp, status); },
clearForm: false
};
this.targetForm.ajaxForm(opts);
}
AjaxFormSwitcher.prototype.checkValidation = function(formData, jqForm, options) {
jqForm.validate();
return jqForm.valid();
}
AjaxFormSwitcher.prototype.processResult = function(afs, xmlHttp, status) {
if (xmlHttp == null) return;
var r = xmlHttp;
var c = r.getResponseHeader("content-type");
if (c.match("json") != null) {
var json = eval("(" + r.responseText + ")");
afs.processJsonRedirect(json);
}
if (c.match("html") != null) {
afs.outputTarget.html(r.responseText);
}
}
AjaxFormSwitcher.prototype.processJsonRedirect = function(data) {
if (data!=null) {
switch (data.command) {
case 'redirect':
window.location.href = data.content;
break;
}
}
}
That little script handles stuff like wiring up the form to do ajax and processing the result (either json command or html which gets displayed in the target). I admittedly suck at writing javascript, so there is probably a much more graceful way to write that.
share|improve this answer
nice post, some useful ideas there. – Damien Aug 6 '09 at 17:58
1
"Mom! Phineas and Ferb are answering questions on Stack Overflow!" – Charlie Flowers Jan 8 '10 at 6:25
1
???. Or who are Phineas and Ferb and who's mom are you talking about? – Wyatt Barnett Jan 11 '10 at 21:40
1
It's a kids' show on Disney (yes, I have a kid), and they have a running gag in which they say, "Yes, yes I do". – Charlie Flowers Jan 12 '10 at 0:24
Disclaimer: I haven't done much ASP.NET MVC Programming.
That said, I've used a number of MVC frameworks in other languages, and have completed quite a few django projects. Over time, I've kinda evolved a pattern to deal with this sort of thing in django, using template inclusion.
Essentially, I create a partial template containing the form (and any validation errors), which gets included in the main template. The view (or Controller, in your case) then selects between these two templates: AJAX requests get the partial template, and regular requests get the full template. In the main template, I use the jQuery form plugin to submit the form via AJAX, and simply replace the current form with the data from the server: if validation fails, I get a form with highlighted fields, and a list of errors on top; if the POST succeeds, the form is replaced with a success message, or whatever is more appropriate for your situation.
I guess in the ASP.NET MVC world, UserControls could be the equivalent of partials? This article shows how to acheive something similar to what I've described (though, the article seems rather dated, and things could've changed).
To answer your second question, I don't think you should be doing a redirect on "page not found" -- one returns a 302, the other should return a 404; one isn't an error condition, the other is. If you lose the status code, your javascript gets more complicated (since you'll have to somehow test the actual returned data, to figure out what happened). These two posts should give you some ideas on how to implement HTTP-friendly error pages. Django does something similar (raise a Http404 exception, return a 404 status code and optionally render a 404.html template).
share|improve this answer
Action Selector Attribute for Ajax and normal requests
Best way of distinguishing between Ajax requests and normal ones is to write a custom action selector attribute (AjaxOnlyAtribute) and provide two action methods that each handles it's own situation.
[HttpPost]
[AjaxOnly]
[ActionName("Add")]
public ActionResult AddAjax(Entity data) { ... }
[HttpPost]
[ActionName("Add")]
public ActionResult AddNormal(Entity data) { ... }
This way you'll avoid code branches and keep code footprint small while also maintaining control over your code. Provide the paired actions only to those actions that need them.
Handling validation errors in Ajax requests
Handling validation errors in Ajax calls (or any other errors basically) can be done using exception action filter. I've written a particular one called ModelStateExceptionAttribute. This is the way it's done:
[HandleModelStateException]
public ActionResult SomeAjaxAction(Data data)
{
if (!this.ModelState.IsValid)
{
throw new ModelStateException(this.ModelState);
}
// usual code from here on
}
And you've probably made your Ajax request using jQuery so such errors gent easily handled by error ajax function:
$.ajax({
type: "POST",
url: "someURL",
success: function(data, status, xhr) {
// handle success
},
error: function(xhr, status, err) {
// handle error
}
});
You can read a detailed blog post about this approach here.
share|improve this answer
We never redirect (why make the user repeatedly press the 'back' button in case they don't understand what needs to be inputted into a particular field?), we display errors on the spot, be it AJAX or not (the location of the 'spot' on the page is entirely up to you, for all AJAX requests we just show a coloured bar top of the page, just like stackoverflow does for first-comers, only ours does not push the rest of the content down).
As for form validation, you can do both server-side and client-side. We always try to show server-side errors on top of the form in a distinctive container, client-side - right next to the field in question on submit. Once the page comes back from the server with server-side validation errors, user will only see the server-side ones initially, so no duplication there.
For data specified in two places, not sure if I understand, since I've never dealt with nHibernate Validation Mappings.
share|improve this answer
He doesn't want to have to define his validation logic twice: he already defines it using his nHibernate validation mapping attributes (applied to class properties) and doesnt want to have to do it again in his JQuery client-side javascript. I assume he wants to be able to create the client-side javascript validation from the validation in his nHibernate mappings. – Merritt Aug 4 '09 at 17:16
Yes, I am aware this support will exist in MVC 2 but until then. – Damien Aug 4 '09 at 17:32
Hmm.. So nHibernate would be considered server-side then? Still, in my understanding, client-side validation should always complement server-side validation to prevent incorrect data from trickling into the DB, mostly because it can be really tedious to fix later, as well as for security purposes. – dalbaeb Aug 4 '09 at 17:33
Yeah, Server side is done. This isn't so much an issue really. – Damien Aug 4 '09 at 20:44
Keep two error views, one for normal (complete) display of the page and the other which just renders the errors in XML or JSON.
Never validate client-side as the user can easily bypass it. For real-time validation, just run an AJAX request to the server for validation, that way you're writing the validation code once.
share|improve this answer
I have used xVal in the past, along with some homegrown reflection-based JS rule generators. The idea is that you define the rule once (in your case, via nHibernate), and an HTML helper reflects over your properties and generates client-side validation code on the fly (using the jQuery.validation plugin). Exactly the right way to have a responsive client-side UI, while still enforcing server-side validation.
Unfortunately, this method doesn't work for AJAX-posted forms.
For AJAX rules, it would be as simple as adding a Errors array to your returned JSON objects. Anywhere you're using AJAX, simply check the length of Errors (that's all ModelState.IsValid does) and display an error. You can use the IsAjaxRequest method to detect an AJAX call:
public ActionResult PostForm(MyModel thing)
{
UpdateModel(thing);
if (this.Request.IsAjaxRequest() == false)
{
return View();
}
else
{
foreach(var error in ModelState.Errors)
{
MyJsonObject.Errors.Add(error.Message);
}
return JsonResult(MyJsonObject);
}
}
share|improve this answer
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Permit
Order – Perciformes
Family – Carangidae
Genus – Trachinotus
Species – Falcatus
Taxonomy
info_permitCarolus Linnaeus originally described the permit in 1758, classifying it within the jack family carangidae. He initially named it Labrus falcatus, but later taxonomists reclassified the permit within the genus Trachinotus, grouping it with similar species such as the Florida pompano (Trachinotus carolinus), the palometa (Trachinotus goodei), and the blackblotch pompano (Trachinotus kennedyi). Though the permit has been confused as all of these fish, Trachinotus falcatus serves as the only valid scientific name for the species. Falcatus, a Latin adjective that translates to “armed with scythes,” appropriately describes the large sickle-shaped dorsal fin that breaches the surface when permit feed.
Geographical Distribution
Permit inhabit the western Atlantic from Massachusetts to southeastern Brazil. They occur throughout the West Indies and the Gulf of Mexico, and less-frequently in Bermuda. The species has been reported in the eastern Atlantic, but does not regularly occur there. The species is most abundant in southern Florida.
Habitat Permit primarily occupy inshore regions such as flats and sandy beaches, and deeper cuts, channels, and holes adjacent to these areas. The substrate of the flats may vary from sand, mud, marl, or sea grass. Permit often swim in water depths less than 2 feet, though due to large body depth, large individuals cannot occupy waters as shallow as other flats species such as bonefish. In deeper waters up to 30 m, permit often congregate around structures such as reefs, jetties, and wrecks where they frequently occur in large schools.
info_distribution_permit
Biology
Distinctive Features
The deeply forked tail and elongated anterior dorsal fin provide the more distinct characteristics of the permit. Looking like long sickles, these fins impart the fish’s species name falcatus. However, one can also identify the permit by its highly laterally compressed body, making the fish appear thin and tall. From a lateral perspective, the permit shape looks round in juveniles, but becomes oblong as the fish ages into an adult. In addition to the long anterior dorsal fin, inserted directly above an elongated anterior anal fin, permit also possess 17-21 soft dorsal rays, and 16-19 soft anal rays.
Coloration
Permit have bright silver sides and bluish-green or brown backs. The belly will sometimes show yellow or an occasional black splotch. The fins appear dark gray or black.
Dentition
Permit have no teeth other than granular teeth that occur on the tongue, designed for crushing mollusks and crustaceans. The teeth are most conspicuous in younger fish.
Size, Age & Growth
Permit reach a maximum length of at least 48 inches (122 cm) and a weight of 79 pounds (36 kg). They grow rapidly until an age of 5 years, at which point growth slows considerably. Permit reach sexual maturity at about 2.3 years for males, and 3.1 years for females. Their size at sexual maturity ranges from 19.1 inches (486 mm) for males and 21.5 inches (547 mm) for females. Permit can attain an age of 23 years, though they probably live longer. Other carangids, such as the Florida pompano (Trachinotus carolinus) and the white trevally (Pseudocaranx dentex), are known to reach ages of 7 and 49 years, respectively.
Food Habits
Permit primarily forage on flats and intertidal areas, entering shallow water on incoming tides from deeper adjacent channels and basins. They usually travel in schools of about ten, but may school in larger numbers; larger permit tend to be more solitary, feeding alone or in pairs. Permit also congregate around wrecks and other deeper-water structures. Like the bonefish, the permit uses its hard mouth to dig into the benthos and root up its prey. These food items usually consist of crustaceans and mollusks, which the permit crushes with its granular teeth and pharyngeal bony plates. However, as opportunistic feeders, permit will eat a variety of animals, including amphipods, copepods, mollusks, polychaetes, fish and insects. Developmentally, permit exhibit planktivorous feeding habits as juveniles, eating copepods, amphipods, mysids, larval shrimp, and fish. As they increase in size, permit begin to feed on benthic prey including mole crabs, coquin clams, flatworms, gastropods, and sessile barnacles. Larger adults feed on gastropods, sea urchins, bivalves, and crabs.
Reproduction
Permit spawning may last all year, but occurs primarily from May through June in the Florida Keys. Spawning peaks during these summer months, with extended spawning seasons occurring outside this main period and a decrease in spawning activity during the winter months. Researchers have found that permit may spawn over natural and artificial reefs or in near shore waters in the middle and lower Florida Keys. Males reach sexual maturity earlier than females, 2.3 versus 3.1 years, and at respective sizes of 19.1 inches (486 mm) and 21.5 inches (547 mm).
Parasites
Permit suffer from many parasites, including the branchial parasite Bicotylophora baeri, the intestinal parasite Serrasentis socialis, and trematodes from the genus Lobatostoma. When cultured together in fish farms, permit may become susceptible to bacteria such as Vibrio anguillarum and the dinoflagellate Amyloodinium ocellatum. The latter infests the skin and gills of the host, eventually causing death through toxins or by interfering with respiration.
Importance to Humans
Permit compose an important commercial fishery along with their close relative the Florida pompano. The permit commercial fishery yielded 10.4 metric tons in 2002, down from 68 metric tons in 2000. Florida landings comprised 100 percent of the catch in 2002. In response to possible over fishing, the Florida Wildlife Commission recently raised minimum size limits and decreased bag limits (see Conservation below). Originally interested in the related Florida pompano, fish farmers have only recently begun to experiment with the mariculture of permit, raising them in large near-shore pens for commercial sale. Sport fishers consider the permit an important game fish, and this fish, in addition to the bonefish and tarpon, supports a large charter boat fishing industry. Many anglers regard the permit as one of the most difficult game fish to catch, and consider a permit caught on fly the highlight of their angling achievements. Many fishing guides and anglers highly esteem the permit and release the fish unharmed. Most mortality attributed to human activity while sport fishing occurs from injuries incurred when being landed, such as “gut hooking” or sharks that take advantage of the hooked fish. Though conscientious anglers attempt to break the line, thereby releasing the permit from restraint when a shark is sighted, sharks occasionally leave the angler with only half a fish. Danger to Humans Permit pose no threat to humans, though ciguatera poisoning may result if eaten.
Conservation Status
After recent surveys in 2003, the Florida Fish and Wildlife Commission changed the Florida size and bag limits for permit, effective January 1, 2004. The FWC raised the minimum size limit from 10 inches to 11 inches for all fishermen (commercial and recreational), and decreased the recreational aggregate bag limit of permit and Florida pompano from 10 fish to 6 fish per person per day. The upper size limit remains at 20 inches, but with a provision that allows anglers to retain one fish over 20 inches in their daily bag limit. These changes are designed to allow more permit to reach sexual maturity and reduce overall landings. Though susceptible to over fishing, the World Conservation Union (IUCN) does not currently list the permit as an endangered or vulnerable species. The IUCN (a global union of states, governmental agencies, and non-governmental organizations in a partnership) assesses the conservation status of species.
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Getting Started with the Actions SDK
This document guides you through deploying a simple Conversation Action that uses SSML. The action asks you to say a number and it speaks it back to you as an ordinal (for example, if you say "1", the action says "first").
Setup your environment
This tutorial assumes you're using Google Cloud Functions, but you can use any web hosting platform to deploy your Conversation Action fulfillment logic.
1. Download and install Node.js
2. Download, install, and initialize the Google Cloud SDK.
3. Download gactions and save it to a directory we'll refer to as <sdk-dir>
4. Go to the Google Cloud Console projects page. Select a project if you already have one you want to use or create a new project to host your fulfillment logic. Make a note of the project ID, because it is used for configuration and deployment.
5. Enable billing for your project. See the console's help center for more information.
6. Create a Storage Bucket for your Google Developer Project. Make a note of the storage bucket name.
Get the sample and client library
To get the the sample and client library, run the following commands:
git clone https://github.com/actions-on-google/actionssdk-say-number-nodejs.git
cd actionssdk-say-number-nodejs
npm install
Deploy fulfillment and preview the action package
Previewing the action package lets you make it available in a staging area that's only available to the account that is logged in to the gactions tool. To do this, you need to deploy fulfillment to a web server and then preview the action package with the gactions CLI.
1. Deploy the Node.js Cloud Function for fulfillment with the following commands:
cd actionssdk-say-number-nodejs
gcloud beta functions deploy sayNumber --trigger-http --stage-bucket <storage_bucket_name>
You can read server logs with
gcloud beta functions logs read
2. Modify the action.json file in the following places:
• agentInfo.projectId - Specify your Google developer project ID
• httpExecution.url - Specify fulfillment URL of the deployed Node.js Cloud Function from step 1
3. Run the following commands from the actionssdk-say-number-nodejs directory:
<sdk-dir>/gactions preview -action_package=action.json -invocation_name="my action"
You should receive a message to authorize gactions to access your account and Google developer project. Follow the onscreen instructions to do so.
4. Preview the action package by running the following command:
<sdk-dir>/gactions simulate
5. Enter talk to my action to invoke the action.
Sample Explanation
The following directory structure describes the important files for this sample:
<sdk-dir>
|- gactions
|- <sample-dir>
|- actions-on-google.js - Node.js client library
|- action.json - The action package that defines your actions
|- index.js - The fulfillment logic for Google Cloud Functions
|- package.json - Node.js metadata
The two most important files are index.js and action.json, the two main parts of this Conversation Action.
• action.json - This action package defines one action with an initialTrigger that declares a "assistant.intent.action.MAIN" intent (required for every action package). This intent is triggered when you say "talk to my action".
When the main intent is triggered, the Google Assistant hands off the user experience to the Conversation Action, and the action welcomes the user.
• index.js - There are two main functions in the fulfillment logic. mainIntent() handles the StandardIntents.MAIN intent being triggered and welcomes the user. The other function, rawInput(), processes the StandardIntents.TEXT intent. If a user says a number, the action speaks it back as an ordinal (for example, if a user says "one", the action speaks back "first"). If the user says "bye", the action quits.
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Skip to Newest Articles Skip to Search About Us Skip to simple menu
Archive for the 'Headaches' Category
Jaw Pain and Chiropractic
Pain in the jaw joint, or temporomandibular joint dysfunction (TMD), is a common complaint for many people. In fact, a 2009 publication reported that TMD affects between 20-30% of the adult population, typically between 20 and 40 years old, and is more common in females than males. TMD is also the second most frequent cause [..]
Tension-Type Headaches – Management Strategies
Headaches (HA) can significantly alter a person’s quality of life. Moreover, they can interfere and sometimes even prevent an individual from performing important activities such as going to work, attending school, or participating in group activities such as sports, music programs, holiday gatherings, and more. The focus of this month’s article is on tension-type headaches [..]
Posture and Headaches
Headaches (HA) play a significant role in a person’s quality of life and are one of the most common complaints that chiropractors see. This comes as no surprise, as one survey reported 16.6% of adults (18 years and older) suffered from migraines or other severe headaches during the last three months of 2011. Another study [..]
Headaches: How Does Chiropractic Help?
Headaches (HA) can be tremendously disabling, forcing sufferers away from work or play into a dark, quiet room to minimize any noise and light that intensifies the pain. According to the National Headache Foundation, there are over 45 million Americans who suffer from chronic, re-occurring headaches, of which 28 million are of the migraine variety. [..]
Chiropractic and Sinus Headaches
Sinus headaches refer to pain in the head, typically in and around the face. Most of us are knowledgeable about two of our four sinuses: the frontal (forehead) and maxillary (our “cheek bones”). The other two sinuses (called ethmoid and sphenoid) are much less understood. Many patients ask chiropractors about sinus problems, as all of [..]
What Kind of Headache Do I Have? (Part 2)
Previously, we discussed three types of headaches: tension headaches (the most common), cluster headaches (a vascular headache—less common, short duration but REALLY painful), and sinus headaches. Migraine headaches were discussed the month before last. In keeping with the theme, ONE more headache type will be discussed: rebound headaches, followed by anti-inflammatory herbal remedies, and finally, [..]
What Kind of Headache Do I Have?
Headaches come in several different sizes, shapes, and colors, so to speak. In fact, if you search “headache classification,” you will find the IHS (International Headache Society) 152-page manual (PDF) lists MANY different types of headaches! Last month, we discussed migraine headaches. This month, we’ll talk about the other headache types. So WHY is this [..]
Is it a Migraine?
There are MANY different types of headaches, of which migraines are a common type. This discussion will concentrate on some unique characteristics that are associated with migraine headaches. This information may help you understand what type of headache you’re having. A unique feature of migraine headaches is that prior to the start of the headache, [..]
Headaches and Chiropractic Care
There are many different types of headaches with a multitude of symptoms including achy, throbbing, nausea, vomiting, dizziness, numbness, blinding, noise, light and/or odor sensitivity, and more. The causes of headaches can include genetics (familial traits, like migraine headaches), stress or tension (probably one of the most common), environmental (allergies, seasonal, bright sunlight, loud noises, [..]
Do Chiropractors Help Patients With Headaches?
This seems like an easy question to answer, doesn’t it? The answer of course being, YES!!! However, there are many people who suffer with headaches who have never been to a chiropractor or have not even ever considered it as a “good option.” So, rather than having me “reassure you” that chiropractic works GREAT for [..]
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My child has fairly straight teeth and our dentist recommended an orthodontic consultation – is it necessary?
My daughter is 12 years old and has straight top teeth. The bottom teeth are only slightly crooked, however at my daughter’s last check-up the dentist recommended that we see a specialist orthodontist to discuss orthodontic treatment. A couple of parents at school are questioning why I am taking her for a consultation. Why would our dentist be sending us to see the orthodontist when my daughter’s teeth look fine as they are?
Orthodontic treatment is about more than just straight teeth and a beautiful smile. Orthodontic treatment aims to give the patient an excellent occlusion (the way the teeth bite together), so that the teeth and jaws function properly (for speech, eating, etc) but also do not wear down prematurely, which can cause many dental problems later in life. Of course, if your child is undertaking orthodontic treatment, you would also expect them to have a wonderful smile at the end of treatment!
Your daughter may have a deep bite, which is a greater than normal vertical overlap of the bottom teeth by the top teeth. This can cause premature wearing down of the teeth or permanent damage to the gum behind the top front teeth. Once the teeth are worn, or the gum is damaged, this is permanent, and may be difficult to remedy with your dentist.
Despite your daughter’s top teeth being straight, they could also be protruded (or more forward than normal) relative to the lower teeth. Statistics show that children and teenagers with protruding top teeth are more likely to experience trauma to these teeth, resulting in pain as well as lengthy and expensive dental treatment to fix the damage. Sometimes teeth can actually be knocked out and lost permanently. It is not always possible to achieve an ideal cosmetic result with remedial dental treatment such as bridges and implants.
Crossbites may also be present toward the back of your daughter’s mouth. This means that one or more of the top teeth are biting on the inside of the lower teeth, instead of the other way around. Not only can crossbites cause wearing down of the teeth, but they are also often associated with facial asymmetry (the appearance of a crooked face).
Even though you mention that your daughter’s bottom teeth are only slightly crooked, crowding (or crossing over of the teeth) tends to become worse with time, often accelerating during the late teens and early twenties. As we age, our lips tend to drop, so that with maturity, we tend to show more lower teeth than upper teeth. So even on purely cosmetic grounds, if the crowding of your daughter’s lower teeth is likely to become worse with time, it would be better to have them straightened now (and then kept straight in the long term with retainers).
Furthermore, teeth tend to move more quickly and efficiently in adolescents than they do in adults, as orthodontic tooth movement relies on bone metabolism. (It is still possible to move teeth in adults, although more slowly). During puberty or the adolescent growth spurt, the body’s metabolism is very fast, which tends to result in optimum tooth movement and treatment.
So although at first glance your daughter’s teeth look fine, leading you or others to think that orthodontic treatment is unnecessary, there may be very important reasons for your daughter to have an orthodontic consultation. This way, she can be equipped with a healthy functioning bite, resulting in teeth that last a lifetime, to go with her beautiful smile.
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🚀Announcing Flightcontrol - Easily Deploy Blitz.js and Next.js to AWS 🚀
Back to Documentation Menu
RPC Specification
Topics
Jump to a Topic
Caution
This is documentation for using Blitz queries and mutations in a non-Blitz app. For example, from a mobile application or some other app.
For using queries and mutations in your app, refer to Using Queries and Using Mutations.
If you want to access your app e.g. as an external API, take a look at API Routes instead.
Request Details
From an HTTP perspective there is no difference between a query and mutation. They are both function calls that get transformed into POST requests.
Warms up the lambda. For example, when a component renders on the client that uses a mutation, it will issue a HEAD request to the mutation endpoint immediately on page load so that there's a much higher chance the actual mutation request will hit a warm lambda.
Must always return HTTP 200 OK
HEAD Responses
HTTP/1.1 200 OK
POST
Executes the query or mutation.
Request body must be JSON with the following keys:
• params (required) - This will be provided as the first argument to the query or mutation
• version (optional) - String containing the version of the built RPC client
Response body is JSON with the following keys:
• result - the exact result returned from the query/mutation function or null if there was an error
• error - null if success, otherwise an object.
Example POST body:
{
"params": {
"where": {
"id": 1
}
}
}
POST Responses
Success
HTTP/1.1 200 OK
Content-Type: application/json
{
"result": {
"name": "Hello World",
"description": "This is awesome :)"
},
"error": null
}
Bad Request (Not JSON)
HTTP/1.1 400 Bad Request
Content-Type: application/json
{
"result": null,
"error": {
"message": "Request body is not valid JSON"
}
}
Bad Request (Missing params)
HTTP/1.1 400 Bad Request
Content-Type: application/json
{
"result": null,
"error": {
"message": "Request body is missing the 'params' key"
}
}
Not Found (Wrong method)
HTTP/1.1 404 Not Found
Versioning
Any time you have client-server communication, there is a very real risk of the client & server code becoming out of sync and causing an error. For example, let's say you have a mutation function whose input is a single object. You decide to change the mutation input to be a single array. You make the change and deploy it. Your server code is now updated, but all your users still have the old client code running in the browser. So it's likely you now have client code sending an object to your mutation which is going to error because it's expecting an array.
Now one solution is to only make changes that are backwards compatible, but sometimes you can't do this.
The other solution is to track the version of your client code and your server code. If there is a version mismatch, you can display a friendly message to the user, telling them they need to reload the webpage to continue.
// *: example of an RPC call shielded by a friendly message to the user
Idea for improving this page? Edit it on GitHub.
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Hedgehog Fact Sheet Ks1 (Expert Answers)
✅ Fact Checked
Updated on January 16, 2023
Michael Colt, Bachelor Computer Science Degree & Computer Engineering.
Written by
Michael Colt, Bachelor Veterinary Medicine & Animal Science.
Ella Williams
Fact Checked by
Ella Williams
Dr. Michael Colt is a highly qualified veterinarian and animal scientist. He has extensive knowledge and experience in the care and treatment of animals, and a deep understanding of the latest scientific research in the field. Dr. Colt is dedicated to promoting the health and well-being of animals, and is committed to providing the highest level of care to his patients. Holds a Bachelors Degree in Veterinary Medicine from Middle Tennessee State University.
⭐ Fun Fact ⭐
Hedgehogs are known for their ability to roll into a tight ball when they feel threatened. This ball shape helps to protect their soft underbelly and makes it difficult for predators to attack them.
Hedgehogs are fascinating creatures that have captured the hearts of people around the world. They are known for their cute, round faces and spiky exteriors, but there is so much more to these animals than meets the eye. In this fact sheet, we will explore the physical characteristics, habitat, diet, reproduction, behavior, and conservation of hedgehogs. Whether you are a student, teacher, or simply someone who loves animals, this fact sheet will provide you with a comprehensive understanding of these amazing creatures. So, let’s dive in and learn all about hedgehogs!
1 Physical Characteristics of Hedgehogs
Physical Characteristics of Hedgehogs:
Hedgehogs are known for their unique and distinctive appearance. Their bodies are compact and round, with a snout that is short and rounded. This body shape is perfectly designed for their lifestyle, as it allows them to easily roll into a tight ball for protection.
In terms of size, hedgehogs can range from 5 to 12 inches in length, with a weight of around 1-2 pounds. Despite their small size, they are surprisingly strong and can move quickly when they need to.
The fur of a hedgehog is soft and dense, with a mix of brown, black, and white hairs. This fur provides insulation and helps to keep the hedgehog warm in colder temperatures.
One of the most distinctive features of a hedgehog is its spines. These spines are actually modified hairs that are stiff and sharp, and they provide the hedgehog with a powerful defense mechanism. When a hedgehog feels threatened, it will roll into a tight ball, presenting its spines to any potential predators. This makes it difficult for predators to attack the hedgehog, as they are unable to get past the spines.
So essentially, the physical characteristics of hedgehogs are perfectly adapted to their environment and lifestyle. From their compact body shape to their soft fur and sharp spines, these animals are truly remarkable creatures that are well-equipped to survive in the wild.
2 Habitat and Distribution
Habitat and Distribution of Hedgehogs:
Hedgehogs are found in a variety of habitats, including forests, grasslands, and deserts. They are most commonly found in Europe, Africa, and Asia, and they have been introduced to New Zealand and some islands in the Pacific.
See also How To Give A Hedgehog A Foot Bath? (FAQ)
In their natural habitat, hedgehogs will burrow into the ground, using their strong front legs to dig into the soil. They will also use leaves, grasses, and other materials to build a nest, which provides them with a safe and warm place to rest.
The geographical range of hedgehogs is quite extensive, and they are found in many different countries around the world. They are particularly abundant in Europe, where they are considered a common species. In Africa, they are found in a variety of habitats, including savannas, forests, and deserts.
In Asia, hedgehogs are found in a number of different countries, including India, Pakistan, and Afghanistan. They are also found in parts of the Middle East, including Iran and Iraq.
So basically, hedgehogs are adaptable creatures that can be found in a variety of habitats and geographical locations around the world. Whether they are burrowing into the ground in Europe or hiding in the deserts of Africa, these amazing animals are well-equipped to survive and thrive in their environments.
3 Diet and Hunting
Diet and Hunting of Hedgehogs:
Hedgehogs are omnivores, which means that they eat both plants and animals. Their diet mainly consists of insects, such as beetles, caterpillars, and slugs, as well as earthworms, snails, and other small invertebrates. They also eat fruits, berries, and nuts, and will occasionally consume small mammals, such as mice and shrews.
When hunting, hedgehogs will use their strong sense of smell to locate their prey. They will then use their sharp teeth and claws to capture and kill their food. Hedgehogs are nocturnal animals, which means that they are most active at night, and they will spend much of their time foraging for food.
In terms of food preferences, hedgehogs will eat a wide variety of foods, and they are not particularly picky eaters. However, they do have a preference for insects, and will often focus their hunting efforts on finding these tasty treats.
See also Do Hedgehogs Grow New Quills (Detailed Response)
In summary, the diet and hunting habits of hedgehogs are perfectly adapted to their lifestyle and environment. Whether they are foraging for insects or hunting small mammals, these amazing animals are well-equipped to find the food they need to survive and thrive.
4 Reproduction and Life Cycle
Reproduction and Life Cycle of Hedgehogs:
Hedgehogs are solitary animals, and they only come together to mate. Mating usually takes place in the spring, and the male will often follow the female for several days, trying to win her over. Once the female is receptive, mating will take place and the female will become pregnant.
The gestation period for a hedgehog is around 35-45 days, and the female will give birth to a litter of 1-7 young. These young, known as hoglets, are born blind and deaf, and they are completely dependent on their mother for survival.
In terms of litter size, hedgehogs will typically give birth to 2-6 young, although litters of up to 7 have been recorded. The young will nurse from their mother for several weeks, and they will begin to venture out on their own once they are strong enough.
As such, the reproduction and life cycle of hedgehogs are fascinating and complex. From the mating habits of the adults to the growth and development of the young, these amazing animals provide us with a glimpse into the incredible world of wildlife.
5 Behaviour and Adaptations
Behaviour and Adaptations of Hedgehogs:
Hedgehogs are solitary animals, and they do not form social bonds with other hedgehogs. They are primarily active at night, and they will spend much of their time foraging for food and building nests.
In terms of defense mechanisms, hedgehogs are well-equipped to protect themselves from predators. When they feel threatened, they will roll into a tight ball, presenting their sharp spines to any potential predators. This makes it difficult for predators to attack the hedgehog, as they are unable to get past the spines.
See also What Are Hedgehogs A Sign Of? (Detailed Response)
In addition to their spines, hedgehogs have a number of other adaptations that help them to survive in the wild. For example, they have strong front legs that are perfectly adapted for digging, and they have a keen sense of smell that helps them to locate their food.
So essentially, the behaviour and adaptations of hedgehogs are perfectly suited to their environment and lifestyle. Whether they are foraging for food, building nests, or defending themselves from predators, these amazing animals are truly remarkable creatures that are well-equipped to survive in the wild.
6 Threats and Conservation
Threats and Conservation of Hedgehogs:
Hedgehog populations are facing a number of threats, including habitat loss, road traffic, and pesticide use. These threats have led to declines in hedgehog populations in many parts of the world, and they are considered a species of concern in many areas.
To address these threats, conservation efforts are underway to protect and conserve hedgehog populations. These efforts include habitat restoration and protection, road safety measures, and the promotion of sustainable agriculture practices.
In addition, a number of organizations are working to raise awareness about the importance of hedgehogs and the threats they face. These organizations are working to educate the public about the importance of hedgehogs, and they are encouraging people to take action to protect these amazing animals.
So essentially, the threats facing hedgehog populations are significant, but there is hope. Through conservation efforts and public education, we can work to protect and conserve these amazing animals for future generations to enjoy.
7 FAQ
What are 5 interesting facts about hedgehogs?
Yes, of course! Here are five interesting facts about hedgehogs:
1. Hedgehogs are omnivores, meaning they eat both plants and animals, with insects being their main food source.
2. Hedgehogs have a unique defense mechanism, they can roll into a tight ball to protect their soft underbelly and present their sharp spines to potential predators.
See also What Does A Hedgehog Eat As A Pet? (FAQ)
3. Hedgehogs have strong front legs that are adapted for digging, and they use these legs to build nests and burrow into the ground.
4. Hedgehogs are nocturnal animals and are most active at night, spending much of their time foraging for food.
5. Hedgehogs have a keen sense of smell, which they use to locate food and navigate their environment.
What are hedgehogs info for kids?
Yes, here is some information about hedgehogs for kids:
1. Hedgehogs are small, spiny mammals that are found in many parts of the world.
2. They have a unique defense mechanism, they can roll into a tight ball to protect their soft underbelly and present their sharp spines to potential predators.
3. Hedgehogs are omnivores, meaning they eat both plants and animals, with insects being their main food source.
4. They are nocturnal animals and are most active at night, spending much of their time foraging for food.
5. Hedgehogs are important to their ecosystem, as they help to control insect populations and spread seeds through their droppings.
This information should give kids a good understanding of hedgehogs and their role in the world.
What is a fact sheet about hedgehogs?
Yes, a fact sheet about hedgehogs is a document that provides information about these small spiny mammals. It typically includes information about their physical characteristics, habitat, diet, behavior, and life cycle, as well as any threats they may face and conservation efforts underway to protect them. Fact sheets are a great way to learn about hedgehogs and their role in the world, and they are often used by educators, scientists, and conservationists to share information about these fascinating animals.
What are three interesting facts about hedgehogs?
Yes, here are three interesting facts about hedgehogs:
1. Hedgehogs are omnivores, meaning they eat both plants and animals, with insects being their main food source.
See also What Do Hedgehog Quills Feel Like? (Fact Checked)
2. They have a unique defense mechanism, they can roll into a tight ball to protect their soft underbelly and present their sharp spines to potential predators.
3. Hedgehogs are nocturnal animals and are most active at night, spending much of their time foraging for food.
8 Conclusion
In conclusion, hedgehogs are fascinating and important animals that play a critical role in our ecosystems. From their habitat and distribution, to their diet and hunting habits, to their reproduction and life cycle, and their behavior and adaptations, hedgehogs are truly remarkable creatures that deserve our attention and protection.
As we look to the future, it is clear that hedgehog populations are facing a number of threats, including habitat loss, road traffic, and pesticide use. However, through conservation efforts and public education, we can work to protect and conserve these amazing animals for future generations to enjoy.
So, let’s take a moment to appreciate the importance of hedgehogs and the role they play in our world. By working together, we can ensure that these amazing animals have a bright and secure future.
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fulltext.study @t Gmail
Dehydrogenation of n-butane over Pd–Ga/Al2O3 catalysts
Paper ID Volume ID Publish Year Pages File Format Full-Text
42225 45916 2010 5 PDF Available
Title
Dehydrogenation of n-butane over Pd–Ga/Al2O3 catalysts
Abstract
Dehydrogenation of n-butane over alumina-supported Pd and Pd–Ga catalysts was studied. Catalysts were prepared by incipient wetness impregnation with a Pd content of 0.66 wt% and atomic Ga/Pd ratios from 0 to ∞, using aqueous solutions of PdCl2 and Ga(NO3)3. Fresh (uncalcined) and calcined catalysts were characterized by X-ray fluorescence spectroscopy (XRF), N2 adsorption, temperature programmed reduction (TPR), CO and H2 chemisorptions and O2/H2 titrations. n-Butane dehydrogenation reaction was carried out at 500 °C, atmospheric pressure and a H2/C4H10 ratio of 10. An increase in the Pd dispersion with increasing Ga content was observed for the fresh catalysts, according to CO chemisorption results. H2 chemisorption and H2/O2 titrations were not reliable methods to determine the Pd dispersion in the fresh catalysts. For the calcined catalysts with low Ga contents, the CO/Pd, H/Pd and O/Pd values were very similar, but different for those with the higher Ga contents. These differences were explained based on the presence of chlorine. Coke deposition produced the deactivation of the catalysts and inhibited hydrogenolysis reactions on Pd, favoring the dehydrogenation selectivity. In the fresh catalysts, the Ga addition caused an effect similar to that produced by coke, reducing the activity and increasing the dehydrogenation selectivity. In the calcined catalysts, the effect of Ga addition on activity was diminished by the calcination treatment. The combination of calcination and high Ga content led to a catalyst with enhanced activity and very high dehydrogenation selectivity.
Graphical abstractDehydrogenation of n-butane over alumina-supported Pd–Ga catalysts was studied. In the reduced catalysts (without a previous calcination), the Ga addition caused an effect similar to that produced by coke, that is, it reduced the activity and increased the dehydrogenation selectivity. The effect of Ga addition on catalytic activity was partially inhibited by the calcination treatment. At higher Ga contents, the calcination led to catalysts with enhanced activities and very high dehydrogenation selectivities.Figure optionsDownload full-size imageDownload high-quality image (174 K)Download as PowerPoint slide
Keywords
Pd–Ga/Al2O3 catalysts; n-C4H10 dehydrogenation; TPR; CO chemisorption; O2/H2 titration
First Page Preview
Dehydrogenation of n-butane over Pd–Ga/Al2O3 catalysts
Publisher
Database: Elsevier - ScienceDirect
Journal: Applied Catalysis A: General - Volume 373, Issues 1–2, 31 January 2010, Pages 66–70
Authors
, , , , , ,
Subjects
Physical Sciences and Engineering Chemical Engineering Catalysis
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The Truth Behind Energy Drinks: Pros and Cons
Do you rely on energy drinks like Reign, Bang, or Monster to push through your afternoon energy slump? Discover the pros and cons behind this habit.
woman-drinking-green-drink
by
Caroline Thomason
— Signos
Dietician + Diabetes Educator (CDCES)
Green checkmark surrounded by green circle.
Updated by
Green checkmark surrounded by green circle.
Science-based and reviewed
Published:
April 23, 2024
March 22, 2024
— Updated:
Table of Contents
The demand for a quick energy fix has led to our universal acceptance and popularity of energy drinks. Loved for their immediate effect on our level of alertness, physical performance, and a surge of energy, these beverages have become a staple for many seeking an immediate pick-me-up. However, it’s important to understand the potential side effects and consequences of energy drinks. By understanding the pros and cons surrounding these beverages, you can make informed choices regarding their consumption and prioritize their health by making the best choice for you personally.
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What Are Energy Drinks?
Drinks like Monster, Red Bull, and 5-Hour Energy shots all provide a burst of energy from caffeine in large doses. The total amount of caffeine can vary from 80 mg to 200 mg or more in a serving, and having more than one serving without knowing your limits may pose some health risks or side effects.
Having an occasional energy drink likely poses no significant issues and can be incorporated into a balanced diet from time to time. However, if you notice a pattern of daily reliance or use them in specific situations, like before a workout, it may be beneficial to reflect on why you depend on these bottled boosts.
Certain energy drink brands might even incorporate nutrients such as vitamins, minerals, and amino acids into their formulas. Although these supplementary ingredients don't fully offset the potential negative impacts of regular energy drink consumption, they offer a slight nutritional boost, helping mitigate some of the drawbacks or side effects. L-theanine, for example, is often added to slow down the absorption of caffeine, providing a more stable stream of energy without the jitters or crashing later.
Natural vs. Synthetic Caffeine
When shopping for energy drinks, it’s important to know that there are two main caffeine sources: natural and synthetic. While both forms can provide the desired stimulating effects, natural sources like coffee and green tea often contain additional compounds, such as antioxidants, that may mitigate some adverse effects of caffeine consumption.
Conversely, synthetic caffeine, commonly found in energy drinks and supplements, lacks these natural buffers and may pose a higher risk of side effects such as jitteriness, anxiety, and increased heart rate, particularly when consumed in excessive amounts. There have recently been recalls from synthetic caffeine for adverse reactions. It is likely safer to choose caffeine sourced from natural sources like coffee beans, tea, and chocolate for now.
When Are Energy Drinks Healthy?
a cup of pineapple juice
Energy drinks may potentially be suitable for individuals who are in good health, have no underlying medical conditions, and consume them in moderation. Here are some scenarios where energy drinks might be considered relatively healthy.1
• Athletic Events: Energy drinks can provide a quick energy boost before intense workouts or competitions, enhancing performance and endurance.
• Shift Work: Those who work long or irregular hours may benefit from the temporary alertness provided by energy drinks to combat fatigue and maintain productivity.
• Demanding Schedules: During heightened academic or professional demands, energy drinks may offer a short-term solution to enhance focus and concentration.
• Night Activities: Whether driving long distances or participating in all-night events, energy drinks can help individuals stay awake and alert when necessary. If you regularly rely on energy drinks to pull an all-nighter, this can quickly become an unhealthy habit.
However, moderation is key. Additionally, individuals should be aware of their personal caffeine sensitivity, any underlying health conditions, and potential interactions with medications before consuming energy drinks. You should consult with a healthcare professional, especially if you have preexisting health concerns.
What Are the Risks of Energy Drinks?
Caffeine in energy drinks offers a quick jolt of energy and a heightened sense of alertness, making them go-to choices for those needing a productivity boost or a pre-workout kick. However, the potent caffeine content can quickly turn from a perk to a problem, leading to jitteriness, disrupted sleep patterns, and even dependence. Overindulgence can be troubling for those with health conditions affected by caffeine, with potential risks including increased heart rate, elevated blood pressure, and gastrointestinal distress. Moderation is key, as navigating the pros and cons of caffeine in energy drinks requires a careful balance between searching for a buzz without the backlash.
Here are some reasons why energy drinks would be considered unhealthy.2
• Sugar: The recommended daily intake for added sugar is 24 g for women and 36 g for men. Many energy drinks not only meet but often exceed these limits. Opting for sugar-free alternatives can help control your added sugar intake effectively.
• Caffeine: With a recommended daily limit of 400 mg, equivalent to about four cups of coffee, it's easy to surpass this threshold with energy drinks. Consuming more than one energy drink in a short span could lead to exceeding this limit.
• Dependence: Relying on energy drinks to sustain energy levels can lead to dependence, making it challenging to respond to moderate caffeine sources like coffee. This reliance may escalate, necessitating higher doses of energy drinks to achieve the desired alertness.
<p class="pro-tip"><strong>Also Related: </strong><a href="period-brain-fog">Period Brain Fog: Your Menstrual Cycles and Mental Health</a>.</p>
Who Should Not Drink Energy Drinks?
Energy drinks account for a significant portion of emergency room trips per year.3 Related symptoms range from insomnia to irregular heart rate to panic attacks. Here are subsets of the population that are typically deterred from regularly using energy drinks.4
• Children and Teenagers: Most energy drinks are inappropriate for children or teens due to the caffeine content. Teenagers are recommended to limit their caffeine consumption to about a cup of coffee daily, or 100 mg.
• Pregnant and Nursing Women: Caffeine in pregnancy has been connected to a host of negative health outcomes for a newborn baby, and the recommended limit is 200 mg per day. However, many energy drinks contain this amount in addition to other ingredients that may be harmful in pregnancy or breastfeeding.
• Physical Symptoms: Excessive consumption of energy drinks can induce symptoms such as anxiety, elevated heart rate, and jitteriness, indicating an overdose of caffeine. Recognizing and addressing these symptoms promptly is crucial for maintaining overall well-being. For folks with anxiety, gastrointestinal stress, or cardiovascular disease, energy drinks might be contraindicated entirely. Talk to your doctor if you have a medical health history before introducing energy drinks.
• Mixing With Alcohol: Mixing energy drinks with alcohol is a common practice among young adults and college students, posing significant health risks. The caffeine in energy drinks can mask the effects of alcohol, leading to increased alcohol consumption, higher rates of drinking and driving, and more alcohol-related injuries. Studies have shown a heightened risk of heart palpitations when combining energy drinks with alcohol.5 Despite regulatory efforts from the FDA, many continue to mix these beverages, emphasizing the need to avoid this combination due to its dangers.
Occasionally, indulging in an energy drink for a quick boost is acceptable, but it's essential to exercise caution to prevent overconsumption. For a healthy choice, drink less than 400 mg of caffeine daily, and look for brands that contain minimal added sugars.
Other Ways to Boost Energy Levels
a guy sleeping while hugging a pillow
Typically, most health professionals are not big advocates for energy drinks. Instead, registered dietitians and healthcare providers encourage individuals to better understand why they feel the need for these beverages and concentrate on lifestyle adjustments that could address those underlying factors.
• Sleep: If you're feeling drained in the afternoon, perhaps you are not getting adequate or quality sleep. No quantity of energy drinks can substitute for a good night's rest. Both quality of sleep and quantity of sleep are important for maintaining good energy throughout the day and may help prevent over-reliance on energy drinks.
• Nutrition: People who rely on energy drinks tend to skip meals for convenience, opting for a quick fix instead. Prioritizing proper nutrition in the form of regular meals and snacks throughout the day can sustain your energy levels, potentially reducing your need for more caffeine.
• Exploring Alternatives: Coffee, tea, or even a bit of chocolate can provide a natural energy boost with added health benefits. Incorporating these options in moderation can be a healthier substitute for energy drinks. Remember that your recommended caffeine intake from all sources is 400 mg per day, according to the Food and Drug Administration (FDA).6
Learn More About Healthy Nutrition With Signos’ Expert Advice
Managing stable energy levels begins with healthy blood sugar levels, which help provide stable brain power and physical energy levels and might impact mental health. The bottom line: Gaining insight into the glycemic index of foods could lead you to take greater control of your energy levels. Understanding the blood sugar-energy connection can improve your health and wellness.
The expert guidance provided by Signos can significantly benefit your health, whether your goal is better energy or simply feeling better. Explore more about nutrition and adopt healthier habits by delving into Signos' blog, or discover if Signos' program is right for you through a brief quiz.
<p class="pro-tip"><strong>Learn More: </strong><a href="low-sugar-drinks">11 Refreshing Low-Sugar Drinks That Won't Spike Blood Sugar</a>.</p>
Get more information about weight loss, glucose monitors, and living a healthier life
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Topics discussed in this article:
References
1. Alsunni AA. Energy Drink Consumption: Beneficial and Adverse Health Effects. Int J Health Sci (Qassim). 2015 Oct;9(4):468-74.
2. Alsunni AA. Energy Drink Consumption: Beneficial and Adverse Health Effects. Int J Health Sci (Qassim). 2015 Oct;9(4):468-74.
3. Mattson ME. Update on Emergency Department Visits Involving Energy Drinks: A Continuing Public Health Concern. 2013 Jan 10. In: The CBHSQ Report. Rockville (MD): Substance Abuse and Mental Health Services Administration (US); 2013–. PMID: 27606410.
4. Alsunni AA. Energy Drink Consumption: Beneficial and Adverse Health Effects. Int J Health Sci (Qassim). 2015 Oct;9(4):468-74.
5. Peacock A, Bruno R, Martin FH. The subjective physiological, psychological, and behavioral risk-taking consequences of alcohol and energy drink co-ingestion. Alcohol Clin Exp Res. 2012 Nov;36(11):2008-15.
6. Spilling the beans: How Much Caffeine is Too Much? FDA. Accessed 3/15/24.
About the author
Caroline Thomason is a dietitian, diabetes educator, and health writer based in Washington, DC.
View Author Bio
Please note: The Signos team is committed to sharing insightful and actionable health articles that are backed by scientific research, supported by expert reviews, and vetted by experienced health editors. The Signos blog is not intended to diagnose, treat, cure or prevent any disease. If you have or suspect you have a medical problem, promptly contact your professional healthcare provider. Read more about our editorial process and content philosophy here.
Interested in learning more about metabolic health and weight management?
Try Signos.
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| 0.573454 |
Non-BRCA1/2 Breast Cancer Susceptibility Genes: A New Frontier with Clinical Consequences for Plastic Surgeons
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Abstract
Summary:
Twenty percent of breast cancer cases may be related to a genetic mutation conferring an increased risk of malignancy. The most common and prominent breast cancer susceptibility genes are BRCA1 and BRCA2, found in nearly 40% of such cases. However, continued interest and investigation of cancer genetics has led to the identification of a myriad of different breast cancer susceptibility genes. Additional genes, each with unique significance and associated characteristics, continue to be recognized. Concurrently, advanced genetic testing, while still controversial, has become more accessible and cost-effective. As oncologic and reconstructive advances continue to be made in prophylactic breast reconstructive surgery, patients may present to plastic surgeons with an increasingly more diverse array of genetic diagnoses to discuss breast reconstruction. It is therefore imperative that plastic surgeons be familiar with these breast cancer susceptibility genes and their clinical implications. We, therefore, aim to review the most common non-BRCA1/2 breast cancer susceptibility genetic mutations in an effort to assist plastic surgeons in counseling and managing this unique patient population. Included in this review are syndromic breast cancer susceptibility genes such as TP53, PTEN, CDH1, and STK11, among others. Nonsyndromic breast cancer susceptibility genes herein reviewed include PALB2, CHEK2, and ataxia telangiectasia mutated gene. With this knowledge, plastic surgeons can play a central role in the diagnosis and comprehensive treatment, including successful breast reconstruction, of all patients carrying genetic mutations conferring increased risk for breast malignancies.
Related Topics
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| 0.875046 |
Pedro Vale Estrela
Learn More
All the proposed IP mobility protocols assume that the mobile nodes always have a mobility-aware IP stack. On the other hand, efficient micro-mobility solutions entail specific topologies and mobile-aware routers, requiring major changes in the existing infra-structures. Major advantages are foreseen if mobility can be supported using the existing legacy(More)
All IP mobility protocols currently proposed by the IETF assume that the mobile nodes always have a mobility-aware IP stack, which is still a scenario that can seldom be found nowadays. Most terminals, including the laptops and PDAs which would most benefit of the mobility support, still use legacy IP stacks, limiting their use to layer 2 mobility within a(More)
In the field of mobility support, several mobility protocols resort to the use of triangulation mechanisms as a means of supporting fast handovers or basic connectivity. In order to reduce the maximum end-to-end delay of the packets, such triangulations can be later removed to enable direct routing of the data packets. However, using a simple update of this(More)
This paper describes a Terminal's Independent Mobility Architecture, TIMIP / sMIP, which supports IP mobility of legacy terminals, and presents a complete performance evaluation to compare the proposed architecture with alternative solutions, via simulation studies in complex scenarios featuring wired mesh topologies and wireless disjoint channels. The NS2(More)
This paper describes a local Route Optimization (RO) scheme applicable to IP micro-mobility protocols, which increases the protocols' efficiency and scalability while maintaining their fast-handovers capabilities. This is done by reducing and decentralizing the chain of agents that participate in the data forwarding to stationary mobile nodes. This generic(More)
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• The VIC Vascular Team
DVT and me...
What is DVT and how does it develop? To be healthy and live and active life, we have to have balanced circulation. DVT is the formation of a thrombis or clot formed of red blood cells and this is when circulation fails in the body. The development of a thrombis can block a vessel or travel within that vessel, this occurs for three main reasons. The first two, stasis of blood, and injury to a blood vessel, can lead to formation of a clot. This is why doctors and nurses insist on people moving after surgery or after a trauma or injury. Thirdly, “thick blood” can occur in people who have cancer or clotting abnormalities and lead to a clot forming. DVT is a clot on the vein side of the circulation. This is where blood comes back to the heart and lungs from the body.
What are the common symptoms of DVT? Are they always obvious?
The most common symptoms are pain and swelling. Both of these can be sudden and severe. It is almost always noticed by the patient. Sometimes during certain illnesses or surgeries the risk of DVT can be higher.
How is DVT diagnosed and treated?
DVT is diagnosed with the use of an ultrasound machine to “see” the blood clot or look for abnormal flow. DVT is treated almost universally with a blood thinner. Massive DVT or PE can be treated by vascular surgeons who can remove or dissolve the clot to improve blood flow.
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Mehryar.com
Alica P. Craig
The Most Efficient Weight Loss Diet: And the winner is…. Weight Watchers newest program, Beyond the Scale, is designed to help individuals eat higher and exercise more, thus helping followers create a healthier life-style that can have a long-lasting effect. In comparison with no therapy, exercise helped people lose a small amount of weight, and when people began to train and cut their energy, they misplaced extra weight than with a diet alone.
And air air pollution is a particularly bad fats increaser: A 2011 examine from the Faculty of Public Well being at Ohio State University found just that: Publicity to positive particulate matter (air pollution) induced insulin resistance, diminished glucose tolerance, and increased irritation, leading researchers to mark long-time period exposure to air pollution as a risk issue for diabetes.
Analysis has shown that along with serving to with sustained weight reduction, exercise can have several different positive results on our lives, including decreased threat for chronic disease (think coronary heart illness, cancer and diabetes), effects of stress and high blood pressure. Much less fats will get your muscle tissues observed, no matter what activity you do.
Properly, there are various methods for shedding those additional kilos, which embrace, cardio and energy teaching, having a correctly-balanced weight reduction plan that features further fruit and veggies and fewer processed food , chopping down on sugar and excessive-glycemic index (GI) meals, and the like. Brief-term dramatic weight reduction is rarely healthy or sustainable over time.
Your finest guess is to stay as intently as you can to complete plant foods. But a brand new research finds people who follow a low-carb diet after slimming down could have an edge at maintaining the pounds off. The current iteration of the diet followed by many dietitians suggests seven to 10 servings of vegatables and fruits each day, and fish once or twice per week, but both amounts are adjustable based on way of life and taste.
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Diet foods that style great and may also help you shed pounds. And when obese adults had been assigned to a regular lowered-calorie weight loss diet or a lacto-ovo vegetarian diet, they misplaced roughly the same quantity of weight over the 18 months researchers from the University of Pittsburgh tracked them, […]
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| 0.519821 |
ISCA Archive Eurospeech 1991
ISCA Archive Eurospeech 1991
Pitch estimation based on a "narrowed" autocorrelation function
N. Dal Degan, M. Fratti
This paper presents experiments and results relative to a pitch determination of speech signals based on the "Narrowed" Autocorrelation Function (NACF), recently introduced by Brown and Puckette [1], It is shown that a reliable pitch estimation can be obtained even in a very noisy environment. An advantage of the NACF is that no preprocessing is needed before performing the pitch detection task. Also, a simple post-processing technique is sufficient to recover from the rare, isolated gross pitch errors that may occur. Finally, the paper proposes a short-cut in the computation of the NACF and it is shown that the number of multiplications can be drastically reduced. Keywords: Pitch, Narrowed autocorrelation
doi: 10.21437/Eurospeech.1991-115
Cite as: Degan, N.D., Fratti, M. (1991) Pitch estimation based on a "narrowed" autocorrelation function. Proc. 2nd European Conference on Speech Communication and Technology (Eurospeech 1991), 457-460, doi: 10.21437/Eurospeech.1991-115
@inproceedings{degan91_eurospeech,
author={N. Dal Degan and M. Fratti},
title={{Pitch estimation based on a "narrowed" autocorrelation function}},
year=1991,
booktitle={Proc. 2nd European Conference on Speech Communication and Technology (Eurospeech 1991)},
pages={457--460},
doi={10.21437/Eurospeech.1991-115}
}
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| 0.622255 |
1. Computer problem? Tech Support Guy is completely free -- paid for by advertisers and donations. Click here to join today! If you're new to Tech Support Guy, we highly recommend that you visit our Guide for New Members.
Windows Media SDK
Discussion in 'Windows XP' started by Rude_Dog, May 10, 2007.
Thread Status:
Not open for further replies.
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1. Rude_Dog
Rude_Dog Thread Starter
Joined:
Dec 18, 2006
Messages:
43
(Compaq 2100 laptop, WinXPhome)
What is Windows Media SDK, and why is it on my computer? I see it pop up on Spybot S&D scans. AFAIK, SDK stands for Software Development Kit. Since I'm not developing any Windows Media software, I figure I don't need it on my computer. Should I remove it? It doesn't show up in the Add/Remove programs list. How do I remove it?
Thanks....
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3. Rude_Dog
Rude_Dog Thread Starter
Joined:
Dec 18, 2006
Messages:
43
Bump!
OK, so from another thread that I started and was closed, we have (looks like MS boilerplate):
I've never installed any Windows SKD items, what's it doing on my computer?
I checked add/remove programs, but Windows Media SDK isn't listed as something to remove. If it's not essential, I'd really rather disable it or remove it.
Any ideas?
Thanks
4. dvk01
dvk01 Derek Moderator Malware Specialist
Joined:
Dec 14, 2002
Messages:
50,243
Just because it appears in a spybot scan doesn't mean you actually have it on your computer
Is it appearing in the list of scanned objects or tasks that spybot does or is it in an actual spybot report
You are not giving a lot of information to be able to help you
5. Rude_Dog
Rude_Dog Thread Starter
Joined:
Dec 18, 2006
Messages:
43
When Spybot scans, one of the things it reports is "usage tracks", which are usually information items stored in the registry about how I've used my computer. In a typical scan, 3 of these items appear under the heading Windows Media SDK.
Perhaps my question should be what's making these registry entries that Spybot seems to indicate are from Windows Media SDK?
I'll post again with the specific registry entries as soon as they show up in another scan.
6. dvk01
dvk01 Derek Moderator Malware Specialist
Joined:
Dec 14, 2002
Messages:
50,243
they are just useage tracks and are nothing to do with Windows Media SDK as described earlier but are spybot's name for the media player useage tracks
just ignore them or delete them as you feel suitable
however it's pointless deleting them as they will come back every time you boot windows
7. Rude_Dog
Rude_Dog Thread Starter
Joined:
Dec 18, 2006
Messages:
43
OK, so the following items showed up in my Spybot scan under the title Windows Media SDK.
Whereas Media Player usage tracks show up in the Spybot scan under the title MS Media Player as follows:
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| 0.721854 |
Take the 2-minute tour ×
Stack Overflow is a question and answer site for professional and enthusiast programmers. It's 100% free, no registration required.
In Git, file contents are stored in blob objects. Tree objects represent a directory. They list the files in the directory, and the blob in which the file contents are stored. They also list subdirectories, and the tree that represents the subdirectory. In this sense, a tree in Git is a recursive data structure: it can contain references to itself.
The alternative would be a non-recursive data structure, where a single tree object would list all tracked files in the repository, and store not only the filename, but the full path relative to the root of the repository.
I admire the elegant way in which Git treats data, so I am convinced there must be a good reason why Git uses the recursive form. I came up with several reasons, all of which I discarded:
• A tree per subdirectory would allow tracking directories. It could enable tracking empty directories. Storing directory metadata is more elegant in recursive form (the non-recursive tree would require two types of entry to track subdirectories).
However, Git does not track empty directories, nor does it store directory metadata.
• Recursive trees allow unchanged subdirectories to share the same tree among commits. This would be more space efficient in a repository with many files, where only a few files change per commit, especially when many subdirectories are used.
However, space efficiency is not a concern in tree objects: packfiles take care of compression, and the non-recursive trees might even be better compressible because they contain less hashes (which are essentially random). Additionally, the data-to-files ratio will generally be large, so the size of the tree objects has little impact on the repository size.
• Recursive trees might allow for more elegant algorithms. Operations on the repository can be handled recursively.
However, (perhaps recursively) building a list of all files in the repository, and then operating on the flat list, is in many cases even easier.
• Recursive trees could allow for easier detection of directory renames.
However, Git uses heuristics to detect renames anyway, so this is of little use.
I am unsure about the differences in performance. It is not obvious to me that one of the two representations would be faster.
My question is: why is the recursive form used? Is there a specific advantage to using the this form?
share|improve this question
2
"Recursive trees allow unchanged subdirectories to share the same tree among commits. This would be more space efficient in a repository with many files, where only a few files change per commit, especially when many subdirectories are used." --- i think you hit the nail on the head in your question. – xero Aug 5 '13 at 18:27
you seem to be assuming that non-recursive is the default, and there's needs to be a reason for recursive. in my experience, if you're storing trees, you automatically choose a recursive structure, because it tends to work better for all the reasons you list. so even if there's no obvious single explanation, that's the "normal" way to do it. – andrew cooke Aug 5 '13 at 19:19
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| 0.63057 |
Java Enum 101 : Explained with easy examples
Java Enum or Enumeration are mainly used to define some set of named constants that fall under the same group. It is a special type of Java class.
What is java enum?
If I were to give a real-world example, you can consider a set of animals, which can have different names but fall under the same category of animals (ex – horse, goat, cow, lion etc). Here different animals can be referred to as enumerators and the “animals” category as an enumerated type.
The main motive of having enums in JAVA is to be able to define our own data types. It was added in Java version 5 and above.
Simplest java enum example
One simple enum example could be an enum that has different sizes under it. This could be useful if you are creating a web application for an e-commerce platform and in similar use cases.
public enum SIZE {
LARGE,
SMALL,
MEDIUM
}
The enum keyword has a special meaning here. It informs the compiler that similar to other classes or interfaces, we are defining java enum as well.
How to assign value to an enum type?
Now how do we use an enum, or refer to one of the values in it? Just create an Enum object with the required enum value.
Size size = Size.SMALL;
In the above case, the size variable is of type Size and has been assigned a small value. We can assign LARGE and MEDIUM similarly.
Also check : What is the main method in java ? Explained with example
How to declare an enum?
Few things to consider while declaring an enum in java :
Example:
If declared outside a class.
enum Size
{
SMALL, MEDIUM, LARGE;
}
public class Test
{
public static void main(String[] args)
{
Size size = Size.MEDIUM;
System.out.println(size);
}
}
//output MEDIUM
If declared within a class.
public class Test
{
enum Size {
SMALL, MEDIUM, LARGE;
}
public static void main(String[] args)
{
Size size = Size.MEDIUM; System.out.println(size);
}
}
//output MEDIUM
• The first line of an enum data type should be the list of comma-separated constants. After that one can define custom methods and other things.
• Convention for declaring an enum is that constants should be always capitalized.
• An enum is by default private static final hence. Therefore values or constants are accessed using enum name dot constant name.
• Since it is implicitly declared as final, child enums cannot be created.
• Enums can be compared by using:
• == Operator
• or .equals() method.
• We can also declare main() method within an enum as we would do in a class.
enum Size
{
MEDIUM, SMALL, LARGE;
// Driver method
public static void main(String[] args)
{
Size size = Size.MEDIUM;
System.out.println(size);
}
}
//output MEDIUM
• Enum in java cannot extend any other class, as it already extends java.lang.Enum. However it can implement interfaces if required.
Methods of enum in java:
• values()– Values method is used to get all the values of an enum.
• ordinal()– This method is used to determine the index of a constant in an enum
• valueOf()– Given a string, it returns the constant equivalent of that string.
• A java enum can have custom methods that can be user-defined.
• When defining enum with custom methods, the constants list must end with a semicolon(;).
• An enum can contain abstract methods too. The constraint with this is all instances of the enum class will have to implement abstract methods.
• Custom constructors can also be defined. Default constructor can be overloaded with parameterized constructors.
• When the enum is called, it will call the constructor for all the respective constants of enum one by one during class loading.
• Custom constructors defined has to be private.
• Since there cannot be enum constructors we cannot invoke constructors directly.
• Enum in Java is type-safe. You cannot assign any other value apart from the values specified in the enum.
Also check : Best class to java decompilers : Try online or download free
Example of user-defined method within an enum :
enum Size
{
SMALL, MEDIUM, LARGE;
private Size()
{
System.out.println("Size constructor initialized : " +
this.toString());
}
public void sizeDetails()
{
System.out.println("Covers size for all age groups!");
}
}
public class Sizes
{
public static void main(String[] args)
{
Size size = Size.SMALL;
System.out.println(size);
size.sizeDetails();
}
}
Output:
Size constructor initialized : SMALL
Size constructor initialized : MEDIUM
Size constructor initialized : LARGE
Covers size for all age groups!
Defining customized enum in java:
All the above examples used the constants as default values for enums but java enum also supports custom values. We will take a look at how to define such enums, with the example where:
• We will have an enum of books
• Each member of the enum will be a programming language book name.
• Each book will have its own author
enum Books
{
JAVA(“AUTHOR 1”), CPP(“AUTHOR 2”), PYTHON(“AUTHOR 3”);
}
Few pre-requisite for having a custom enum with values are:
• Enum should have a member variable to save value of constant.
• Enum should have a parametrized constructor initializing the created member variable.
• It should have a getter exposing member variable value.
Example :
public enum JavaWorkshop {
JAVA("Author 1"), CPP("Author 2"), PYTHON("Author 3");
private String programmingLangAuthor;
private JavaWorkshop(String s) {
programmingLangAuthor = s;
}
public String getProgrammingLang() {
return programmingLangAuthor;
}
}
Then using it in a class.
public class JavaWorkshopExample {
public static void main(String[] args) {
System.out.println("Get author of book 'JAVA': "
+ JavaWorkshop.JAVA.getProgrammingLang());
}
}
Output:
Get author of book 'JAVA': Author 1
Few Faqs related to Enum in java :
When to use java enum?
Enums are useful and can be used only when the use case satisfies certain criteria.
• If all the possible constant values are known beforehand.
• If a class has to use one value at a time from a list of constants.
• The same set of constants are to be used at multiple places within a class or in multiple classes.
Also check : What is a queue in Java ? Explained with Examples
Is enum immutable?
Enum in java is tightly immutable. They expose certain methods which can be used to get all the values of an enum. This array of values can be used to create a new enum with some additional values as a workaround.
Can enum in java be null?
Yes, a java enum can be null. When a field of type enum is created in a class but not initialized, by default it is null.
Why java enum?
Because they enable reusability and hence better code readability. And them being immutable add an extra layer of reliability for the constants defined within code.
Let’s Discuss
If you liked the tutorial, feel free to share it on social media. If you have any doubts or something is missing from the post, leave them in the comments below.
Leave a Comment
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__label__pos
| 0.996089 |
February 1, 2023
Sulfate, Sulfite, Sulfide
One of the positive trends in the agriculture industry is the push toward natural farming and food production methods that are environmentally friendly. They promote healthy soils, and limit the potential for exposure to toxic chemicals in the food chain. The wine industry has been on the leading edge with organic farming methods, as well as low impact wine making in the cellar.
One substance that leads to mass confusion is sulfur. I see this frequently in dealing with consumers, so this month I will address the role of sulfur in wine making, and why in limited amounts it’s considered a winemaker’s best friend.
Sulfate, sulfite, and sulfide are the three common basic oxidation states of the sulfur molecule, or in chemical terms the SO42-, SO32-, and S2- ions respectively. Microbes in wine are able to convert sulfate to sulfite, then sulfide through metabolic pathways. Sulfate, along with thiosulfate (S2O32-), are common sources of sulfur ions in amino acid biosynthesis and occur naturally in grapes in the form of potassium sulfate and potassium thiosulfate.
Sulfite occurs naturally in grapes in small amounts, but is limited in concentration by the vine as its elemental form SO2(aq) disrupts cell membranes. Some amount of sulfite gets produced by yeast during fermentation through sulfate conversion.
Sulfides are a group of organosulfur compounds called thioethers, which are similar to ether, except the oxygen molecule is replaced by sulfur. Compounds such as hydrogen sulfide (H2S) and dimethyl sulfide (DMS) are a part of this group. They’re known for being the stinky sulfur compounds. They’re present in all wines, and in trace amounts add to complexity, but are offensive if too abundant.
Understanding Sulfites
It’s the sulfite form which is the sticking point of the sulfur argument, as that’s what winemakers use as an additive to their wine. To understand how sulfite works, one must understand how it behaves across the pH spectrum.
Wines typically range from a pH of 3.0 to 4.0, with averages around 3.3 to 3.6 for whites and 3.4 to 3.8 for reds. Sulfite has three basic ion forms in solution: the sulfite ion (SO32-), bisulfite (HSO3-), and sulfur dioxide (SO2(aq)), also known as the molecular or elemental form. The sulfite ion is the common form above pH values of 7.18, and is very limited below pH values of 5.0. Bisulfite is the major form between pH values of 1.86 and 7.18, and is the predominant form in wine. Molecular SO2 is the second most common species at wine pH and its concentration as a percentage in solution will range roughly from 5.56% at a pH of 3.0 to 0.585% at a pH of 4.0.
Sulfites are added to wine for two basic reasons: to control unwanted microbial growth, and to preserve the wine from oxidation. Molecular SO2 is the form responsible for antimicrobial activity due to its ability to disrupt cell membranes, effectively killing single celled organisms. Wines that are more acidic with lower pH values require less SO2 to maintain effective kill levels for spoilage microbes. The trend towards over-ripe, full and fruity wine, which tends to have higher pH values, requires more SO2 in the winery.
The bisulfite form is primarily responsible for anti-oxidation. Bisulfite in and of itself is actually a very poor antioxidant. Its binding rates with oxygen are far too slow compared with other compounds present in wine that often react with oxygen enzymatically. What bisulfite does do is bind exceptionally well to the byproducts of those oxidation reactions, essentially sequestering them from sensory impact. Bisulfite also binds with many other compounds, particularly sugar. As a result, a larger total SO2 addition is required at higher pH values, or when wines are sweeter, to have enough free SO2 available in the wine versus what binds to other molecules.
The moral of the science is, if you want less SO2, drink drier and more acidic wines. The push for no SO2 winemaking is not one I’m in favor of, because when used judiciously it keeps wine fresh from oxidation and spoilage. SO2 became a hot button issue as people were told it was an allergen. Molecular SO2 when concentrated is a strong irritant that bothers the mucus membranes, but very few people are truly allergic to it. Those who are usually are severe asthmatics. The truth is, the human body produces far more SO2 per day than is present in an average bottle of wine, and it’s a common food preservative time tested in wine making since ancient Rome.
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__label__pos
| 0.809221 |
Anterior cruciate ligament xenografts
- CrossCart, Inc.
The invention provides an article of manufacture comprising a substantially non-immunogenic ligament or tendon xenograft for implantation into humans. The invention further provides a method for preparing a ligament xenograft by removing at least a portion of an ligament from a non-human animal to provide a xenograft; washing the xenograft in saline and alcohol; subjecting the xenograft to at least one treatment selected from the group consisting of exposure to ultraviolet radiation, immersion in alcohol, ozonation, freeze/thaw cycling, and optionally chemical crosslinking. In addition to or in lieu of the above treatments, the methods include a cellular disruption treatment and glycosidase digestion of carbohydrate moieties of the xenograft followed by treatment of carbohydrate moieties of the xenograft with capping molecules. The invention also provides articles of manufacture produced by one or more of the above-identified methods of the invention. The invention further provides a ligament xenograft for implantation into a human including a portion of a ligament from a non-human animal, wherein the portion includes extracellular components and substantially only dead cells. The extracellular components and dead cells have substantially no surface &agr;-galactosyl moieties and have capping molecules linked to at least a portion of surface carbohydrate moieties. Each of the xenografts of the invention is substantially non-immunogenic and has substantially the same mechanical properties as the respective native ligament.
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Description
RELATED APPLICATIONS
This application is a divisional of application U.S. Ser. No. 09/036,172, filed Mar. 6, 1998, now U.S. Pat. No. 6,110,206 which is a continuation-in-part of the U.S. Ser. No. 08/529,199, filed Sep. 15, 1995, now U.S. Pat. No. 5,902,338 the disclosures of which are incorporated herein by reference.
FIELD OF THE INVENTION
The present invention relates to the field of surgical repair of injuries of the anterior cruciate ligament in the human knee using a substantially immunologically compatible ligament or tendon from a non-human animal to replace the damaged human anterior cruciate ligament.
BACKGROUND OF THE INVENTION
The anterior cruciate ligament of the knee (hereinafter the ACL) functions to resist anterior displacement of the tibia from the femur at all flexion positions. The ACL also resists hyperextension and contributes to rotational stability of the fully extended knee during internal and external tibial rotation. The ACL may play a role in proprioception. The ACL is made up of connective tissue structures composed of cells, water, collagen, proteoglycans, fibronectin, elastin, and other glycoproteins. Cyril Frank , M. D. et al., Normal Ligament: Structure, Function, and Composition. Injury and Repair of the Musculoskeletal Soft Tissues, 2:45-101. Structurally, the ACL attaches to a depression in the front of the intercondyloid eminence of the tibia extending postero-superiorly to the medial wall of the lateral femoral condyle.
Partial or complete tears of the ACL are very common, comprising about 30,000 outpatient procedures in the U.S. each year. The preferred treatment of the torn ACL is ligament reconstruction, using a bone-ligament-bone autograft. Cruciate ligament reconstruction has the advantage of immediate stability and a potential for immediate vigorous rehabilitation. However, the disadvantages to ACL reconstruction are significant: for example, normal anatomy is disrupted when the patellar tendon or hamstring tendons are used for the reconstruction; placement of intraarticular hardware is required for ligament fixation; and anterior knee pain frequently occurs. Moreover, recent reviews of cruciate ligament reconstruction indicate an increased risk of degenerative arthritis with intraarticular ACL reconstruction in large groups of patients.
A second method of treating ACL injuries, referred to as “primary repair”, involves suturing the torn structure back into place. Primary ACL repair has the potential advantages of a limited arthroscopic approach, minimal disruption of normal anatomy, and an out-patient procedure under a local anesthetic. The potential disadvantage of primary cruciate ligament repair is the perception that over the long term ACL repairs do not provide stability in a sufficient number of patients, and that subsequent reconstruction may be required at a later date. The success rate of anterior cruciate ligament repair has generally hovered in the 60% to 70% range.
Much of the structure and many of the properties of original tissues may be retained in transplants through use of xenograft or heterograft materials, that is, tissue from a different species than the graft recipient. For example, tendons or ligaments from cows or other animals are covered with a synthetic mesh and transplanted into a heterologous host in U.S. Pat. No. 4,400,833. Flat tissues such as pig pericardia are also disclosed as being suitable for heterologous transplantation in U.S. Pat. No. 4,400,833. Bovine peritoneum fabricated into a biomaterial suitable for prosthetic heart valves, vascular grafts, burn and other wound dressings is disclosed in U.S. Pat. No. 4,755,593. Bovine, ovine, or porcine blood vessel xenografts are disclosed in WO 84/03036. However, none of these disclosures describe the use of a xenograft for ACL replacement.
Once implanted in an individual, a xenograft provokes immunogenic reactions such as chronic and hyperacute rejection of the xenograft. The term “chronic rejection”, as used herein, refers to an immunological reaction in an individual against a xenograft being implanted into the individual. Typically, chronic rejection is mediated by the interaction of IgG natural antibodies in the serum of the individual receiving the xenograft and carbohydrate moieties expressed on cells, and/or extracellular components. For example, transplantation of ligament or tendon xenografts from nonprimate mammals (e.g., porcine or bovine origin) into humans is primarily prevented by the interaction between the IgG natural anti-Gal antibody present in the serum of humans with the carbohydrate structure Gal&agr;1-3Gal&bgr;1-4G1cNAc-R (&agr;-galactosyl or &agr;-gal epitope) expressed in the xenograft. K. R. Stone et al., Porcine and bovine cartilage transplants in cynomolgus monkey: I. A model for chronic xenograft rejection, 63 Transplantation 640-645 (1997); U. Galili et al., Porcine and bovine cartilage transplants in cynomolgus monkey: II. Changes in anti-Gal response during chronic rejection, 63 Transplantation 646-651 (1997). In chronic rejection, the immune system typically responds within one to two weeks of implantation of the xenograft.
In contrast with “chronic rejection”, “hyper acute rejection” as used herein, refers to the immunological reaction in an individual against a xenograft being implanted into the individual, where the rejection is typically mediated by the interaction of IgM natural antibodies in the serum of the individual receiving the xenograft and carbohydrate moieties expressed on cells. This interaction activates the complement system causing lysis of the vascular bed and stoppage of blood flow in the receiving individual within minutes to two to three hours.
The term “extracellular components”, as used herein, refers to extracellular water, collagen, proteoglycans, fibronectin, elastin, and other glycoproteins present in the ligament or tendon.
Xenograft materials may be chemically treated to reduce immunogenicity prior to implantation into a recipient. For example, glutaraldehyde is used to cross-link or “tan” xenograft tissue in order to reduce its antigenicity, as described in detail in U.S. Pat. No. 4,755,593. Other agents such as aliphatic and aromatic diamine compounds may provide additional crosslinking through the side chain carboxyl groups of aspartic and glutamic acid residues of the collagen polypeptide. Glutaraldehyde and diamine tanning also increases the stability of the xenograft tissue.
Xenograft tissues may also be subjected to various physical treatments in preparation for implantation. For example, U.S. Pat. No. 4,755,593 discloses subjecting xenograft tissue to mechanical strain by stretching to produce a thinner and stiffer biomaterial for grafting. Tissue for allograft transplantation is commonly cryopreserved to optimize cell viability during storage, as disclosed, for example, in U.S. Pat. Nos. 5,071,741; 5,131,850; 5,160,313; and 5,171,660. 5,071,741 discloses that freezing tissues causes mechanical injuries to cells therein because of extracellular or intracellular ice crystal formation and osmotic dehydration.
SUMMARY OF THE INVENTION
The present invention is directed against the chronic rejection of xenografts for implantation into humans. Accordingly, the ligament xenograft produced in accordance with the method of the invention is substantially non-immunogenic, while generally maintaining the mechanical properties of a corresponding portion of a native ligament. While the ligament may undergo some shrinkage during processing, a ligament xenograft prepared in accordance with the invention will have the general appearance of a native ligament. The ligament xenograft may also be cut into segments, each of which may be implanted into the knee of a recipient as set forth below.
As described herein, the term “ligament” also includes tendons.
As described herein, the term “xenograft” is synonymous with the term “heterograft” and refers to a graft transferred from an animal of one species to one of another species. Stedman's Medical Dictionary, Williams & Wilkins, Baltimore, Md. (1995).
As described herein, the term “xenogeneic”, as in xenogeneic graft ligament, etc., refers to a graft, ligament, etc., transferred from an animal of one species to one of another species. Id.
The methods of the invention, include, alone or in combination, treatment with radiation, one or more cycles of freezing and thawing, treatment with a chemical cross-linking agent, treatment with alcohol or ozonation. In addition to or in lieu of these methods, the methods of the invention include a cellular disruption treatment and glycosidase digestion of carbohydrate moieties of the xenograft followed by treatment of carbohydrate moieties of the xenograft with capping molecules. After one or more of the above-described processing steps, the methods of the invention provide a xenograft having substantially the same mechanical properties as a corresponding portion of a native ligament.
As described herein, the term “cellular disruption” as in, for example, cellular disruption treatment, refers to a treatment for killing cells.
As described herein, the term “capping molecule(s)”, refers to molecule(s) which link with carbohydrate chains such that the xenograft is no longer recognized as foreign by the subject's immune system.
In one embodiment, the invention provides an article of manufacture comprising a substantially non-immunogenic ligament xenograft for implantation into a human.
In another embodiment, the invention provides a method of preparing a ligament xenograft for implantation into a human, which includes removing at least a portion of a ligament from a non-human animal to provide a xenograft; washing the xenograft in water and alcohol; and subjecting the xenograft to at least one treatment selected from the group consisting of exposure to ultraviolet radiation, immersion in alcohol, ozonation, and freeze/thaw cycling, whereby the xenograft has substantially the same mechanical properties as a corresponding portion of a native ligament
As described herein, the term “portion”, as in, for example, a portion of ligament or a portion of surface carbohydrate moieties, refers to all or less than all of the respective ligament or surface carbohydrate moieties.
In still another embodiment, the invention provides a method of preparing a ligament xenograft for implantation into a human, which includes removing at least a portion of a ligament from a non-human animal to provide a xenograft; washing the xenograft in water and alcohol; subjecting the xenograft to a cellular disruption treatment; digesting the xenograft with a glycosidase to remove substantially first surface carbohydrate moieties from the xenograft; and treating second surface carbohydrate moieties on the xenograft with capping molecules to cap at least a portion of the second surface carbohydrate moieties, whereby the xenograft is substantially non-immunogenic and has substantially the same mechanical properties as a corresponding portion of a native ligament.
As described herein, the terms “to cap” or “capping”, refer to linking a capping molecule, such as a carbohydrate unit, to an end of a carbohydrate chain, as in, for example, covalently linking a carbohydrate unit to surface carbohydrate moieties on the xenograft.
In further embodiments, the invention provides articles of manufacture including substantially non-immunogenic ligament xenografts for implantation into humans produced by one or more of the above-identified methods of the invention.
In still a further embodiment, the invention provides a ligament xenograft for implantation into a human which includes a portion of a ligament from a non-human animal, wherein the portion includes extracellular components and substantially only dead cells, the extracellular components and the dead cells having substantially no surface &agr;-galactosyl moieties and having capping molecules linked to at least a portion of surface carbohydrate moieties. The xenograft is substantially non-immunogenic and has substantially the same mechanical properties as a corresponding portion of a native ligament.
DETAILED DESCRIPTION OF THE PREFERRED EMBODIMENTS
The present invention is directed against the chronic rejection of xenografts for implantation into humans. Accordingly, the ligament xenograft produced in accordance with the method of the invention is substantially non-immunogenic, while generally maintaining the mechanical properties of a native ligament. While the ligament may undergo some shrinkage during processing, a ligament xenograft prepared in accordance with the invention will have the general appearance of a native ligament. The ligament xenograft may also be cut into segments, each of which may be implanted into the knee of a recipient as set forth below.
The invention provides, in one embodiment, a method for preparing or processing a xenogeneic ligament for engraftment into humans. The ligament may be harvested from any non-human animal to prepare the xenograft of the invention. Ligaments from transgenic non-human animals or from genetically altered non-human animals may also be used as xenografts in accordance with the present invention. Preferably, bovine joints serve as sources of the ligament used to prepare the xenografts. More preferably, immature joints from immature animals are the sources of the ligament, since the tissue of younger animals may be inherently more elastic and engraftable than that of older animals. Most preferably, the age of the source animal is between six and eighteen months at time of slaughter. Additionally, the patellar tendon, the anterior or posterior cruciate ligaments, the Achilles tendon, or the hamstring tendons may be harvested from the animal source and used as a donor ligament.
In the first step of the method of the invention, an intact ligament is removed from the knee of a non-human animal. The joint which serves as the source of the ligament should be collected from freshly killed animals and preferably immediately placed in a suitable sterile isotonic or other tissue preserving solution. Harvesting of the joints should occur as soon as possible after slaughter of the animal and preferably should be performed in the cold, i.e., in the approximate range of about 5° C. to about 20° C., to minimize enzymatic degradation of the ligament tissue.
The ligaments are harvested from the joints in the cold, under strict sterile technique. The joint is opened by standard surgical technique. Preferably, the ligament is harvested with a block of bone attached to one or both ends, although in some forms of the invention the ligament alone is harvested. In one form of the invention, a block of bone representing a substantially cylindrical plug of approximately 9-10 mm in diameter by 20-40 mm in length may be left attached to the ligament. The ligament is carefully identified and dissected free of adhering tissue, thereby forming the xenograft.
The xenograft is then washed in about ten volumes of sterile cold water to remove residual blood proteins and water soluble materials. The xenograft is then immersed in alcohol at room temperature for about five minutes, to sterilize the tissue and to remove non-collagenous materials.
After alcohol immersion, the xenograft may be directly implanted into a knee. Alternatively, the xenograft may be subjected to at least one of the following treatments: radiation treatment, treatment with alcohol or ozonation, one or more cycles of freezing and thawing, and treatment with a chemical cross-linking agent. When more than one of these treatments is applied to the xenograft, the treatments may occur in any order.
In one embodiment of the method of the invention, the xenograft may be treated by exposure to ultraviolet radiation for about fifteen minutes or gamma radiation in an amount of about 0.5 to 3 MegaRad.
In another embodiment, the xenograft may be treated by again being placed in an alcohol solution. Any alcohol solution may be used to perform this treatment. Preferably, the xenograft is placed in a 70% solution of isopropanol at room temperature.
In still another embodiment, the xenograft may be subjected to ozonation.
In a further embodiment of the method of the invention, the xenograft may be treated by freeze/thaw cycling. For example, the xenograft may be frozen using any method of freezing, so long as the xenograft is completely frozen, i.e., no interior warm spots remain which contain unfrozen tissue. Preferably, the xenograft is dipped into liquid nitrogen for about five minutes to perform this step of the method. More preferably, the xenograft is frozen slowly by placing it in a freezer. In the next step of the freeze/thaw cycling treatment, the xenograft is thawed by immersion in an isotonic saline bath at room temperature (about 25° C.) for about ten minutes. No external heat or radiation source is used, in order to minimize fiber degradation.
In yet a further embodiment, the xenograft may optionally be exposed to a chemical agent to tan or crosslink the proteins within the extracellular proteins, to further diminish or reduce the immunogenic determinants present in the xenograft. Any tanning or crosslinking agent may be used for this treatment, and more than one crosslinking step may be performed or more than one crosslinking agent may be used in order to ensure complete crosslinking and thus optimally reduce the immunogenicity of the xenograft. For example, aldehydes such as glutaraldehyde, formaldehyde, adipic dialdehyde, and the like, may be used to crosslink the extracellular collagen in accordance with the method of the invention. Other suitable crosslinking agents include aliphatic and aromatic diamines, carbodiimides, diisocyanates, and the like. When glutaraldehyde is used as the crosslinking agent, for example, the xenograft may be placed in a buffered solution containing about 0.05 to about 5.0% glutaraldehyde and having a pH of about 7.4. Any suitable buffer may be used, such as phosphate buffered saline or trishydroxymethylaminomethane, and the like, so long as it is possible to maintain control over the pH of the solution for the duration of the crosslinking reaction, which may be from one to fourteen days, and preferably from three to five days. Alternatively, the xenograft can be exposed to a crosslinking agent in a vapor form, including, but not limited to, a vaporized aldehyde crosslinking agent, such as, for example, vaporized formaldehyde. The vaporized crosslinking agent can have a concentration and a pH and the xenograft can be exposed to the vaporized crosslinking agent for a period of time suitable to permit the crosslinking reaction to occur. For example, the xenograft can be exposed to vaporized crosslinking agent having a concentration of about 0.05 to about 5.0% and a pH of about 7.4, for a period of time which can be from one to fourteen days, and preferably from three to five days. Exposure to vaporized crosslinking agent can result in reduced residual chemicals in the xenograft from the crosslinking agent exposure. The crosslinking reaction should continue until the immunogenic determinants are substantially removed from the xenogeneic tissue, but the reaction should be terminated prior to significant alterations of the mechanical properties of the xenograft. When diamines are also used as crosslinking agents, the glutaraldehyde crosslinking should occur after the diamine crosslinking, so that any unreacted diamines are capped. After the crosslinking reactions have proceeded to completion as described above, the xenograft should be rinsed to remove residual chemicals, and 0.01-0.05 M glycine may be added to cap any unreacted aldehyde groups which remain.
In addition to or in lieu of the above treatments, the xenograft can be subjected to a cellular disruption treatment to kill the xenograft's cells, which precedes or follows digestion of the xenograft with glycosidases to remove surface carbohydrate moieties from the xenograft. The glycosidase digestion in turn can be followed by linkage with capping molecules to cap surface N-acetyllactosamine ends of carbohydrate chains of the xenograft.
In an embodiment of this method of the invention, the xenograft is subjected to a cellular disruption treatment to kill the cells of the ligament prior to in vitro digestion of the xenograft with glycosidases. Typically after surface carbohydrate moieties have been removed from nucleated cells and extracellular components, nucleated cells, i.e., living cells reexpress the surface carbohydrate moieties. Reexpression of antigenic moieties of a xenograft can provoke continued immunogenic rejection of the xenograft. In contrast, non-nucleated, i.e., dead cells, are unable to reexpress surface carbohydrate moieties. Removal of antigenic surface carbohydrate moieties from the non-nucleated cells and extracellular components of a xenograft substantially permanently eliminates antigenic surface carbohydrate moieties as a source of immunogenic rejection of the xenograft.
Accordingly, in the above-identified embodiment, the xenograft of the present invention is subjected to freeze/thaw cycling as discussed above to disrupt, i.e., to kill the cells of the ligament. Alternatively, the xenograft of the present invention is treated with gamma radiation having an amount of 0.2 MegaRad up to about 3 MegaRad. Such radiation kills the ligament cells and sterilizes the xenograft. Once killed, the ligament cells are no longer able to reexpress antigenic surface carbohydrate moieties such &agr;-gal epitopes which are factors in the immunogenic rejection of the transplanted xenografts.
Either before or after the ligament cells are killed, the xenograft is subjected to in vitro digestion of the xenograft with glycosidases, and specifically galactosidases, such as &agr;-galactosidase, to enzymatically eliminate antigenic surface carbohydrate moieties. In particular, &agr;-gal epitopes are eliminated by enzymatic treatment with &agr;-galactosidases, as shown in the following reaction: Gal α1 - 3 Gal β1 - 4 GlcNAc - R α - gal epitope ⟶ α - galactosidase Gal β1 - 4 GlcNAc - R + Gal N - acetyllactosamine
The N-acetyllactosamine residues are epitopes that are normally expressed on human and mammalian cells and thus are not immunogenic. The in vitro digestion of the xenograft with glycosidases is accomplished by various methods. For example, the xenograft can be soaked or incubated in a buffer solution containing glycosidase. In addition, the xenograft can be pierced to increase permeability, as further described below. Alternatively, a buffer solution containing the glycosidase can be forced under pressure into the xenograft via a pulsatile lavage process.
Elimination of the &agr;-gal epitopes from the xenograft diminishes the immune response against the xenograft. The &agr;-gal epitope is expressed in nonprimate mammal and in New World monkeys (monkeys of South America) as 1×106−35×106 epitopes per cell, as well as on macromolecules such as proteoglycans of the extracellular components. U. Galili et al., Man, apes, and Old World monkeys differ from other mammals in the expression of &agr;-galactosyl epitopes on nucleated cells, 263 J. Biol. Chem. 17755 (1988). This epitope is absent in Old World primates (monkeys of Asia and Africa and apes) and humans, however. Id. Anti-Gal is produced in humans and primates as a result of an immune response to &agr;-gal epitope carbohydrate structures on gastrointestinal bacteria. U. Galili et al., Interaction between human natural anti-&agr;-galactosyl immunoglobulin G and bacteria of the human flora, 56 Infect. Immun. 1730 (1988); R. M. Hamadeh et al., Human natural anti-Gal IgG regulates alternative complement pathway activation on bacterial surfaces, 89 J. Clin. Invest. 1223 (1992). Since nonprimate mammals produce &agr;-gal epitopes, xenotransplantation of xenografts from these mammals into primates results in rejection because of primate anti-Gal binding to these epitopes on the xenograft. The binding results in the destruction of the xenograft by complement fixation and by antibody dependent cell cytotoxicity. U. Galili et al., Interaction of the natural anti-Gal antibody with &agr;-galactosyl epitopes: A major obstacle for xenotransplantation in humans, 14 Immunology Today 480 (1993); M. Sandrin et al., Anti-pig IgM antibodies in human serum react predominantly with Gal&agr;1-3Gal epitopes, 90 Proc. Natl. Acad. Sci. USA 11391 (1993); H. Good et al., Identification of carbohydrate structures which bind human anti-porcine antibodies: implications for discordant grafting in man. 24 Transplant. Proc. 559 (1992); B. H. Collins et al., Cardiac xenografts between primate species provide evidence for the importance of the &agr;-galactosyl determinant in hyperacute rejection, 154 J. Immunol. 5500 (1995). Furthermore, xenotransplantation results in major activation of the immune system to produce increased amounts of high affinity anti-Gal. Accordingly, the substantial elimination of &agr;-gal epitopes from cells and from extracellular components of the xenograft, and the prevention of reexpression of &agr;-gal epitopes can diminish the immune response against the xenograft associated with anti-Gal antibody binding with &agr;-gal epitopes.
Following treatment with glycosidase, the remaining carbohydrate chains (e.g., glycosaminoglycans) of the xenograft are optionally treated with capping molecules to cap at least a portion of the remaining carbohydrate chains. Treatment with capping molecules is applicable to both glycosidase-treated and non-glycosidase-treated xenografts, however. For example, xenografts from knock out animals which may lack &agr;-gal epitopes may be treated with capping molecules to cap carbohydrate moieties on the xenograft, thereby reducing the xenograft's immunogenicity. Examples of capping molecules used in the present invention include fucosyl and n-acetyl glucosamine.
Prior to treatment, the outer surface of the xenograft may optionally be pierced to increase permeability to agents used to render the xenograft substantially non-immunogenic. A sterile surgical needle such as an 18 gauge needle may be used to perform this piercing step, or, alternatively a comb-like apparatus containing a plurality of needles may be used. The piercing may be performed with various patterns, and with various pierce-to-pierce spacings, in order to establish a desired access to the interior of the xenograft. Piercing may also be performed with a laser. In one form of the invention, one or more straight lines of punctures about three millimeters apart are established circumferentially in the outer surface of the xenograft.
Prior to implantation, the ligament xenograft of the invention may be treated with limited digestion by proteolytic enzymes such as ficin or trypsin to increase tissue flexibility or coated with anticalcification agents, antithrombotic coatings, antibiotics, growth factors, or other drugs which may enhance the incorporation of the xenograft into the recipient knee joint. The ligament xenograft of the invention may be further sterilized using known methods, for example, with additional glutaraldehyde or formaldehyde treatment, ethylene oxide sterilization, propylene oxide sterilization, or the like. The xenograft may be stored frozen until required for use.
The ligament xenograft of the invention, or a segment thereof, may be implanted into a damaged human knee joint by those of skill in the art using known arthroscopic surgical techniques. Specific instruments for performing arthroscopic techniques are known to those of skill in the art, which ensure accurate and reproducible placement of ligament implants. Initially, complete diagnostic arthroscopy of the knee joint is accomplished using known methods. The irreparably damaged ligament is removed with a surgical shaver. The anatomic insertion sites for the ligament are identified and drilled to accommodate a bone plug. The size of the bone plug can be about 9-10 mm in width by about 9-10 mm in depth by about 20-40 mm in length. The xenogeneic ligament is brought through the drill holes and affixed with interference screws. Routine closure is performed.
This invention is further illustrated by the following Examples which should not be construed as limiting. The contents of all references and published patents and patent applications cited throughout the application are hereby incorporated by reference.
EXAMPLE 1
Assessment of Primate Response to Implanted Bovine Ligament Treated with &agr;-Galactosidase
In this example, bovine ligament implants are treated with &agr;-galactosidase to eliminate &agr;-galactosyl epitopes, the implants are transplanted into cynomolgus monkeys, and the primate response to the ligament implants is assessed.
A bovine stifle joint is sterilely prepared and a ligament with a block of bone attached to one or both ends is removed in the cold, under strict sterile technique. A block of bone representing a substantially cylindrical plug of approximately 9 mm in diameter by 40 mm in length is left attached to the ligament. The ligament is carefully identified and dissected free of adhering tissue, thereby forming the xenograft.
The xenograft is then washed for at least five minutes with an alcohol, such as ethanol or isopropanol, to remove synovial fluid and lipid soluble contaminants.
The ligament specimen is frozen at a temperature of about −70° C. to disrupt, that, is to kill, the ligament specimen's cells.
Each ligament specimen is cut into two portions. The first portion is immersed in a buffer solution containing &agr;-galactosidase at a predetermined concentration. The ligament specimen is allowed to incubate in the buffer solution for a predetermined time period at a predetermined temperature. The second ligament portion is incubated under similar conditions as the first ligament portion in a buffer solution in the absence of &agr;-galactosidase and serves as the control.
At the end of the incubation, the first portion is washed under conditions which allow the enzyme to diffuse out. Assays are performed to confirm the complete removal of the &agr;-gal epitopes.
With the animals under general inhalation anesthesia, the anatomic insertion sites for the xenogeneic ligament are identified and drilled to accommodate a substantially 9 mm in diameter by 40 mm in length bone plug. The xenogeneic ligament is brought through the drill holes and affixed with interference screws. The procedure is performed under sterile surgical technique, and the wounds are closed with 3-0 vicryl or a suitable equivalent known to those of ordinary skill in the art. The animals are permitted unrestricted cage activity and monitored for any sign of discomfort, swelling, infection, or rejection. Blood samples (e.g., 2 ml) are drawn periodically (e.g., every two weeks) for monitoring of antibodies.
The occurrence of an immune response against the xenograft is assessed by determining anti-Gal and non-anti-Gal anti-cartilage antibodies (i.e., antibodies binding to cartilage antigens other than the &agr;-gal epitopes) in serum samples from the transplanted monkeys. At least two ml blood samples are drawn from the transplanted monkeys on the day of implant surgery and at periodic (e.g., two week) intervals post-transplantation. The blood samples are centrifuged and the serum samples are frozen and evaluated for the anti-Gal and other non-anti-Gal anti-cartilage antibody activity.
Anti-Gal activity is determined in the serum samples in ELISA with &agr;-gal-BSA as solid phase antigen, according to methods known in the prior art, such as, for example, the methods described in Galili et al., Porcine and bovine cartilage transplants in cynomolgus monkey: II. Changes in anti-Gal response during chronic rejection, 63 Transplantation 645-651 (1997).
Assays are conducted to determine whether &agr;-galactosidase treated xenografts induce the formation of anti-ligament antibodies. For measuring anti-ligament antibody activity, an ELISA assay is performed according to methods known in the prior art, such as, for example, the methods described in K. R. Stone et al., Porcine and bovine cartilage transplants in cynomolgus monkey: I. A model for chronic xenograft rejection, 63 Transplantation 640-645 (1997).
The ligament xenograft is optionally explanted at one to two months post-transplantation, sectioned and stained for histological evaluation of inflammatory infiltrates. Post-transplantation changes in anti-Gal and other anti-ligament antibody activities are correlated with the inflammatory histologic characteristics (i.e., granulocytes or mononuclear cell infiltrates) within the explanted ligament, one to two months post-transplantation, using methods known in the art, as, for example, the methods described in K. R. Stone et al., Porcine and bovine cartilage transplants in cynomolgus monkey: I. A model for chronic xenograft rejection, 63 Transplantation 640-645 (1997).
Where the xenogeneic ligament is explanted, the ligament xenograft is aseptically harvested. At the time of the xenograft removal, joint fluid, if present in amounts sufficient to aspirate, is collected from the stifle joints for possible immunologic testing if the gross and histopathologic evaluation of the transplants indicate good performance of the transplanted ligament.
A portion of the implant and surrounding tissue is frozen in an embedding medium for frozen tissue specimens in embedding molds for immunohistochemistry evaluation according to the methods known in the prior art. “TISSUE-TEK®” O.C.T. compound which includes 10.24% w/w polyvinyl alcohol, 4.26% w/w polyethylene glycol, and 86.60% w/w nonreactive ingredients, and is manufactured by Sakura FinTek, Torrence, Calif., is a non-limiting example of a possible embedding medium for use with the present invention. Other embedding mediums known to those of ordinary skill in the art may also be used. The remaining implant and surrounding tissue is collected in 10% neutral buffered formalin for histopathologic examination.
EXAMPLE 2
Assessment of Primate Response to Implanted Ligament Treated with &agr;-Galactosidase, Fucosyl and Fucosyltransferase
In this example, bovine ligament implants are treated with &agr;-galactosidase to eliminate &agr;-gal epitopes, as described in Example 1. The implants are further treated with fucosyl and fucosyl transferase to cap carbohydrate chains with fucosyl. Fucosyltransferase facilitates the transfer of fucosyl to the xenograft. The fucosyl links to and thus caps the carbohydrate chains. Capping with fucosyl interferes with the ability of the subject's immune system to recognize the xenograft as foreign. The implants are transplanted into cynomolgus monkeys, and the primate response to the cartilage implants is assessed.
Bovine ligament stifle joints are prepared as described in Example 1 including the &agr;-galactosidase treatment. Prior to implantation into the monkeys, however, the implants are further treated with a predetermined amount of fucosyl and fucosyltransferase, at specified concentrations for a predetermined time and at a predetermined temperature, to cap carbohydrate chains with fucosyl. For example, the sample is immersed in a buffer solution at predetermined concentrations of fucosyl and fucosyl transferase and is incubated for a predetermined time period at a predetermined temperature.
Other molecules, such as n-acetyl glucosamine in combination with the corresponding glycosyltransferase, can also be used for capping the carbohydrate chains of the implants.
Subsequently, the samples are washed to remove the enzyme and implanted into the monkeys, and the occurrence of an immune response against the xenograft is assessed as described above in Example 1.
Those of skill in the art will recognize that the invention may be embodied in other specific forms without departing from the spirit or essential characteristics thereof. The presently described embodiments are therefore to be considered in all respects as illustrative and not restrictive, the scope of the invention being indicated by the appended claims rather than by the foregoing description, and all variations of the invention which are encompassed within the meaning and range of equivalency of the claims are therefor intended to be embraced therein.
Claims
1. A method of preparing a ligament xenograft for implantation into a human, which comprises
a. removing at least a portion of a ligament from a non-human animal to provide a xenograft;
b. washing the xenograft in water and alcohol;
c. subjecting the xenograft to a cellular disruption treatment; and
d. treating a plurality of surface carbohydrate moieties on the xenograft with a plurality of capping molecules to cap at least a portion of the surface carbohydrate moieties,
whereby the xenograft is substantially non-immunogenic and has substantially the same mechanical properties as a corresponding portion of a native ligament.
2. The method of claim 1, wherein at least a portion of the capping molecules are a plurality of fucosyl molecules.
3. The method of claim 1, wherein at least a portion of the capping molecules are a plurality of n-acetyl glucosamine molecules.
4. The method of claim 1, further comprising, after step c, digesting the xenograft with glycosidase to remove substantially a plurality of surface carbohydrate moieties from the xenograft.
5. The method of claim 4, wherein the glycosidase is a galactosidase.
6. The method of claim 5, wherein the galactosidase is an &agr;-galactosidase.
7. The method of claim 1, wherein the cellular disruption treatment comprises freeze/thaw cycling.
8. The method of claim 1, wherein the cellular disruption treatment comprises exposure to gamma radiation.
9. The method of claim 1, wherein the removing step comprises removing with the portion a first block of bone attached to a first end of the portion.
10. The method of claim 9, wherein the removing step comprises removing with the portion a second block of bone affixed to a second end of the portion opposite the first end.
11. The method of claim 1 further comprising the step of following step c, exposing the xenograft to a crosslinking agent in a vapor form.
12. An article of manufacture comprising a substantially non-immunogenic ligament xenograft for implantation in to a human, produced by
a. removing at least a portion of a ligament from a non-human animal to provide a xenograft;
b. washing the xenograft in water and alcohol;
c. subjecting the xenograft to a cellular disruption treatment; and
d. treating a plurality of surface carbohydrate moieties on the xenograft with a plurality of capping molecules to cap at least a portion of the surface carbohydrate moieties,
whereby the xenograft is substantially non-immunogenic and has substantially the same mechanical properties as a corresponding portion of a native ligament.
13. The article of manufacture of claim 12, wherein at least a portion of the capping molecules are a plurality of fucosyl molecules.
14. The article of manufacture of claim 12, wherein at least a portion of the capping molecules are a plurality of n-acetyl glucosamine molecules.
15. The article of manufacture of claim 12, further comprising the step of following step c, digesting the xenograft with a glycosidase to remove substantially a plurality of surface carbohydrate moieties from the xenograft.
16. The article of manufacture of claim 15, wherein the glycosidase is a galactosidase.
17. The article of manufacture of claim 16, wherein the galactosidase is an &agr;-galactosidase.
18. The article of manufacture of claim 12, wherein the cellular disruption treatment comprises freeze/thaw cycling.
19. The article of manufacture of claim 12, wherein the cellular disruption treatment comprises exposure to gamma radiation.
20. The article of manufacture of claim 12, wherein the removing step comprises removing with the portion a first block of bone attached to a first end of the portion.
21. The article of manufacture of claim 12, wherein the removing step comprises removing with the portion a second block of bone affixed to a second end of the portion opposite the first end.
22. The article of manufacture of claim 12 further comprising the step of following step c, exposing the xenograft to a crosslinking agent in a vapor form.
23. A ligament xenograft for implantation into a human comprising
a portion of a ligament from a non-human animal, wherein the portion includes a plurality of extracellular components and a plurality of substantially only dead cells, the extracellular components and the dead cells having a plurality of capping molecules linked to at least a portion of a plurality of surface carbohydrate moieties on the xenograft,
whereby the portion of the ligament is substantially non-immunogenic and has substantially the same mechanical properties as a corresponding portion of a native ligament.
24. The ligament xenograft of claim 23, wherein at least a portion of the capping molecules are a plurality of fucosyl molecules.
25. The ligament xenograft of claim 23, wherein at least a portion of the capping molecules are a plurality of n-acetyl glucosamine molecules.
26. The ligament xenograft of claim 23, wherein the xenograft has substantially no surface &agr;-galactosyl moieties.
27. The ligament xenograft of claim 23, wherein the portion of the ligament has a first block of bone attached to a first end thereof.
28. The ligament xenograft of claim 23, wherein the portion of the ligament has a second block of bone affixed to a second end thereof opposite the first end.
Referenced Cited
U.S. Patent Documents
4034418 July 12, 1977 Jackson et al.
4344193 August 17, 1982 Kenny
4400833 August 30, 1983 Kurland
4502161 March 5, 1985 Wall
4597266 July 1, 1986 Entrekin
4609627 September 2, 1986 Goldstein
4627853 December 9, 1986 Campbell et al.
4642120 February 10, 1987 Nevo et al.
4678470 July 7, 1987 Nashef et al.
4755593 July 5, 1988 Lauren
4776853 October 11, 1988 Klement et al.
4789663 December 6, 1988 Wallace et al.
4801299 January 31, 1989 Brendel et al.
4846835 July 11, 1989 Grande
4880429 November 14, 1989 Stone
4932973 June 12, 1990 Gendler
5007934 April 16, 1991 Stone
5067962 November 26, 1991 Campbell et al.
5071741 December 10, 1991 Brockbank
5078744 January 7, 1992 Chvapil
5092894 March 3, 1992 Kenny
5116374 May 26, 1992 Stone
5131850 July 21, 1992 Brockbank
5158574 October 27, 1992 Stone
5160313 November 3, 1992 Carpenter et al.
5171273 December 15, 1992 Silver et al.
5171322 December 15, 1992 Kenny
5171660 December 15, 1992 Carpenter et al.
5192312 March 9, 1993 Orton
5216126 June 1, 1993 Cox et al.
5306304 April 26, 1994 Gendler
5306311 April 26, 1994 Stone et al.
5358525 October 25, 1994 Fox et al.
5507810 April 16, 1996 Prewett et al.
5613982 March 25, 1997 Goldstein
5944755 August 31, 1999 Stone
Foreign Patent Documents
WO 84/03036 August 1984 WO
WO 95/26740 October 1995 WO
WO 95/28412 October 1995 WO
WO 95/33828 December 1995 WO
Other references
• Rodrigo et al., Clinical Orthopedics and Related Research, 134:342-349 (1978).
• Sengupta et al., The Journal of Bone and Joint Surgery, 56B:167-177 (1974).
• Webber et al., Journal of Orthopedic Research, 3:36-42 (1985).
• Rubak et al., Acta Orthop. Scand, 53:181-186 (1982).
• Engkvist, Ove, Scand. J. Plast. Reconstr. Surg., 13:361-369 (1982).
• Collins et al., Xenotransplantation, Characterization of Porcine Endothelial Cell Determinants Recognized by Human Natural Antibodies, 1:36-46 (1994).
• Satake et al., Xenotransplantation, Limited Specificity of Xenoantibodies In Diabetic Patients Transplanted With Fetal Porcine Islet Cell Clusters. Main Antibody Reactivity Against &agr;-linked Galactose-Containing Epitopes, 1:89-101 (1994).
• LaVecchio et al., Transplantation, Enzymatic Removal of Alpha-Galactosyl Epitopes From Porcine Endothelial Cells Diminishes The Cytotoxic Effect of Natural Antibodies, 60:841-847.
• Stone et al., Arthroscopy: The Journal of Arthroscopic and Related Surgery, 9:234-237 (1993).
• Cotterell et al., Transplantation, The Humoral Immune Response in Humans Following Cross-Perfusion of Porcine Organs, 60:861-868 (1995).
• Galili, Immunology Today, 14:480-482 (1993).
• Elves et al., An Investigation Into The Immunogenicity Of Various Components of Osteoarticular Grafts, The British Journal of Experimental Pathology, 55:344-351 (1974).
Patent History
Patent number: 6402783
Type: Grant
Filed: Jun 28, 2000
Date of Patent: Jun 11, 2002
Assignee: CrossCart, Inc. (San Francisco, CA)
Inventor: Kevin R. Stone (Mill Valley, CA)
Primary Examiner: Corrine McDermott
Assistant Examiner: Choon P. Koh
Attorney, Agent or Law Firm: McDermott, Will & Emery
Application Number: 09/605,222
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| 0.879752 |
Digital signal processing
From Wikipedia, the free encyclopedia
Jump to: navigation, search
Digital signal processing (DSP) is concerned with the representation of the signals by a sequence of numbers or symbols and the processing of these signals. Digital signal processing and analog signal processing are subfields of signal processing. DSP includes subfields like: audio and speech signal processing, sonar and radar signal processing, sensor array processing, spectral estimation, statistical signal processing, digital image processing, signal processing for communications, biomedical signal processing, seismic data processing, etc.
Since the goal of DSP is usually to measure or filter continuous real-world analog signals, the first step is usually to convert the signal from an analog to a digital form, by using an analog to digital converter. Often, the required output signal is another analog output signal, which requires a digital to analog converter. Even if this process is more complex than analog processing and has a discrete value range, the stability of digital signal processing thanks to error detection and correction and being less vulnerable to noise makes it advantageous over analog signal processing for many, though not all, applications. [1]
DSP algorithms have long been run on standard computers, on specialized processors called digital signal processors (DSPs), or on purpose-built hardware such as application-specific integrated circuit (ASICs). Today there are additional technologies used for digital signal processing including more powerful general purpose microprocessors, field-programmable gate arrays (FPGAs), digital signal controllers (mostly for industrial apps such as motor control), and stream processors, among others.[2]
Contents
[edit] DSP domains
In DSP, engineers usually study digital signals in one of the following domains: time domain (one-dimensional signals), spatial domain (multidimensional signals), frequency domain, autocorrelation domain, and wavelet domains. They choose the domain in which to process a signal by making an informed guess (or by trying different possibilities) as to which domain best represents the essential characteristics of the signal. A sequence of samples from a measuring device produces a time or spatial domain representation, whereas a discrete Fourier transform produces the frequency domain information, that is the frequency spectrum. Autocorrelation is defined as the cross-correlation of the signal with itself over varying intervals of time or space.qq
[edit] Signal sampling
With the increasing use of computers the usage of and need for digital signal processing has increased. In order to use an analog signal on a computer it must be digitized with an analog to digital converter (ADC). Sampling is usually carried out in two stages, discretization and quantization. In the discretization stage, the space of signals is partitioned into equivalence classes and quantization is carried out by replacing the signal with representative signal of the corresponding equivalence class. In the quantization stage the representative signal values are approximated by values from a finite set.
The Nyquist–Shannon sampling theorem states that a signal can be exactly reconstructed from its samples if the sampling frequency is greater than twice the highest frequency of the signal. In practice, the sampling frequency is often significantly more than twice the required bandwidth.
A digital to analog converter (DAC) is used to convert the digital signal back to analog. The use of a digital computer is a key ingredient in digital control systems.
[edit] Time and space domains
The most common processing approach in the time or space domain is enhancement of the input signal through a method called filtering. Filtering generally consists of some transformation of a number of surrounding samples around the current sample of the input or output signal. There are various ways to characterize filters; for example:
• A "linear" filter is a linear transformation of input samples; other filters are "non-linear." Linear filters satisfy the superposition condition, i.e. if an input is a weighted linear combination of different signals, the output is an equally weighted linear combination of the corresponding output signals.
• A "causal" filter uses only previous samples of the input or output signals; while a "non-causal" filter uses future input samples. A non-causal filter can usually be changed into a causal filter by adding a delay to it.
• A "time-invariant" filter has constant properties over time; other filters such as adaptive filters change in time.
• Some filters are "stable", others are "unstable". A stable filter produces an output that converges to a constant value with time, or remains bounded within a finite interval. An unstable filter can produce an output that grows without bounds, with bounded or even zero input.
• A "finite impulse response" (FIR) filter uses only the input signal, while an "infinite impulse response" filter (IIR) uses both the input signal and previous samples of the output signal. FIR filters are always stable, while IIR filters may be unstable.
Most filters can be described in Z-domain (a superset of the frequency domain) by their transfer functions. A filter may also be described as a difference equation, a collection of zeroes and poles or, if it is an FIR filter, an impulse response or step response. The output of an FIR filter to any given input may be calculated by convolving the input signal with the impulse response. Filters can also be represented by block diagrams which can then be used to derive a sample processing algorithm to implement the filter using hardware instructions.
[edit] Frequency domain
Signals are converted from time or space domain to the frequency domain usually through the Fourier transform. The Fourier transform converts the signal information to a magnitude and phase component of each frequency. Often the Fourier transform is converted to the power spectrum, which is the magnitude of each frequency component squared.
The most common purpose for analysis of signals in the frequency domain is analysis of signal properties. The engineer can study the spectrum to determine which frequencies are present in the input signal and which are missing.
Filtering, particularly in non realtime work can also be achieved by converting to the frequency domain, applying the filter and then converting back to the time domain. This is a fast, O(n log n) operation, and can give essentially any filter shape including excellent approximations to brickwall filters.
There are some commonly used frequency domain transformations. For example, the cepstrum converts a signal to the frequency domain through Fourier transform, takes the logarithm, then applies another Fourier transform. This emphasizes the frequency components with smaller magnitude while retaining the order of magnitudes of frequency components.
Frequency domain analysis is also called spectrum- or spectral analysis.
[edit] Applications
The main applications of DSP are audio signal processing, audio compression, digital image processing, video compression, speech processing, speech recognition, digital communications, RADAR, SONAR, seismology, and biomedicine. Specific examples are speech compression and transmission in digital mobile phones, room matching equalization of sound in Hifi and sound reinforcement applications, weather forecasting, economic forecasting, seismic data processing, analysis and control of industrial processes, computer-generated animations in movies, medical imaging such as CAT scans and MRI, MP3 compression, image manipulation, high fidelity loudspeaker crossovers and equalization, and audio effects for use with electric guitar amplifiers.
[edit] Implementation
Digital signal processing is often implemented using specialised microprocessors such as the DSP56000, the TMS320, or the SHARC. These often process data using fixed-point arithmetic, although some versions are available which use floating point arithmetic and are more powerful. For faster applications FPGAs[3] might be used. Beginning in 2007, multicore implementations of DSPs have started to emerge from companies including Freescale and startup Stream Processors, Inc. For faster applications with vast usage, ASICs might be designed specifically. For slow applications, a traditional slower processor such as a microcontroller may be adequate.
[edit] Techniques
[edit] Related fields
[edit] References
[edit] Further reading
• Alan V. Oppenheim, Ronald W. Schafer, John R. Buck : Discrete-Time Signal Processing, Prentice Hall, ISBN 0-13-754920-2
• Boaz Porat: A Course in Digital Signal Processing, Wiley, ISBN 0471149616
• Richard G. Lyons: Understanding Digital Signal Processing, Prentice Hall, ISBN 0-13-108989-7
• Jonathan Yaakov Stein, Digital Signal Processing, a Computer Science Perspective, Wiley, ISBN 0-471-29546-9
• Sen M. Kuo, Woon-Seng Gan: Digital Signal Processors: Architectures, Implementations, and Applications, Prentice Hall, ISBN 0-13-035214-4
• Bernard Mulgrew, Peter Grant, John Thompson: Digital Signal Processing - Concepts and Applications, Palgrave Macmillan, ISBN 0-333-96356-3
• Steven W. Smith: Digital Signal Processing - A Practical Guide for Engineers and Scientists, Newnes, ISBN 0-7506-7444-X
• Paul A. Lynn, Wolfgang Fuerst: Introductory Digital Signal Processing with Computer Applications, John Wiley & Sons, ISBN 0-471-97984-8
• James D. Broesch: Digital Signal Processing Demystified, Newnes, ISBN 1-878707-16-7
• John G. Proakis, Dimitris Manolakis: Digital Signal Processing - Principles, Algorithms and Applications, Pearson, ISBN 0-13-394289-9
• Hari Krishna Garg: Digital Signal Processing Algorithms, CRC Press, ISBN 0-8493-7178-3
• P. Gaydecki: Foundations Of Digital Signal Processing: Theory, Algorithms And Hardware Design, Institution of Electrical Engineers, ISBN 0-85296-431-5
• Paul M. Embree, Damon Danieli: C++ Algorithms for Digital Signal Processing, Prentice Hall, ISBN 0-13-179144-3
• Anthony Zaknich: Neural Networks for Intelligent Signal Processing, World Scientific Pub Co Inc, ISBN 981-238-305-0
• Vijay Madisetti, Douglas B. Williams: The Digital Signal Processing Handbook, CRC Press, ISBN 0-8493-8572-5
• Stergios Stergiopoulos: Advanced Signal Processing Handbook: Theory and Implementation for Radar, Sonar, and Medical Imaging Real-Time Systems, CRC Press, ISBN 0-8493-3691-0
• Joyce Van De Vegte: Fundamentals of Digital Signal Processing, Prentice Hall, ISBN 0-13-016077-6
• Ashfaq Khan: Digital Signal Processing Fundamentals, Charles River Media, ISBN 1-58450-281-9
• Jonathan M. Blackledge, Martin Turner: Digital Signal Processing: Mathematical and Computational Methods, Software Development and Applications, Horwood Publishing, ISBN 1-898563-48-9
• Bimal Krishna, K. Y. Lin, Hari C. Krishna: Computational Number Theory & Digital Signal Processing, CRC Press, ISBN 0-8493-7177-5
• Doug Smith: Digital Signal Processing Technology: Essentials of the Communications Revolution, American Radio Relay League, ISBN 0-87259-819-5
• Henrique S. Malvar: Signal Processing with Lapped Transforms, Artech House Publishers, ISBN 0-89006-467-9
• Charles A. Schuler: Digital Signal Processing: A Hands-On Approach, McGraw-Hill, ISBN 0-07-829744-3
• James H. McClellan, Ronald W. Schafer, Mark A. Yoder: Signal Processing First, Prentice Hall, ISBN 0-13-090999-8
• Artur Krukowski, Izzet Kale: DSP System Design: Complexity Reduced Iir Filter Implementation for Practical Applications, Kluwer Academic Publishers, ISBN 1-4020-7558-8
• John G. Proakis: A Self-Study Guide for Digital Signal Processing, Prentice Hall, ISBN 0-13-143239-7
Personal tools
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__label__pos
| 0.81031 |
I have a table with sort on my site. When i load the page the table starts with the 2nd columns sorted from A to Z. I want to change these default settings, but I cant find them in the code. Can you help me please?
This is the table.js file i embed:
Code:
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d[0]=d[0].split('').reverse().join('').replace(/(\d{3})(?=\d)/g,'$1,').split('').reverse().join('');
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function cp(f,c){
var g,h; f=g=f.v.toLowerCase(); c=h=c.v.toLowerCase();
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/*
* File: chromatable.js
* Version: 1.3.0
* CVS: $Id$
* Description: Make a "sticky" header at the top of the table, so it stays put while the table scrolls
* Author: Zachary Siswick
* Created: Thursday 19 November 2009 8:53pm
* Language: Javascript
*
*/
(function($){
$.chromatable = {
// Default options
defaults: {
//specify a pixel dimension, auto, or 100%
width: "900px",
height: "300px",
scrolling: "yes"
}
};
$.fn.chromatable = function(options){
// Extend default options
var options = $.extend({}, $.chromatable.defaults, options);
return this.each(function(){
// Add jQuery methods to the element
var $this = $(this);
var $uniqueID = $(this).attr("ID") + ("wrapper");
//Add dimentsions from user or default parameters to the DOM elements
$(this).css('width', options.width).addClass("_scrolling");
$(this).wrap('<div class="scrolling_outer"><div id="'+$uniqueID+'" class="scrolling_inner"></div></div>');
$(".scrolling_outer").css({'position':'relative'});
$("#"+$uniqueID).css(
{'border':'1px solid #CCCCCC',
'overflow-x':'hidden',
'overflow-y':'auto',
'padding-right':'17px'
});
$("#"+$uniqueID).css('height', options.height);
$("#"+$uniqueID).css('width', options.width);
// clone an exact copy of the scrolling table and add to DOM before the original table
// replace old class with new to differentiate between the two
$(this).before($(this).clone().attr("id", "").addClass("_thead").css(
{'width' : 'auto',
'display' : 'block',
'position':'absolute',
'border':'none',
'border-bottom':'1px solid #CCC',
'top':'1px'
}));
// remove all children within the cloned table after the thead element
$('._thead').children('tbody').remove();
$(this).each(function( $this ){
// if the width is auto, we need to remove padding-right on scrolling container
if (options.width == "100%" || options.width == "auto") {
$("#"+$uniqueID).css({'padding-right':'0px'});
}
if (options.scrolling == "no") {
$("#"+$uniqueID).before('<a href="#" class="expander" style="width:100%;">Expand table</a>');
$("#"+$uniqueID).css({'padding-right':'0px'});
$(".expander").each(
function(int){
$(this).attr("ID", int);
$( this ).bind ("click",function(){
$("#"+$uniqueID).css({'height':'auto'});
$("#"+$uniqueID+" ._thead").remove();
$(this).remove();
});
});
//this is dependant on the jQuery resizable UI plugin
$("#"+$uniqueID).resizable({ handles: 's' }).css("overflow-y", "hidden");
}
});
// Get a relative reference to the "sticky header"
$curr = $this.prev();
// Copy the cell widths across from the original table
$("thead:eq(0)>tr th",this).each( function (i) {
$("thead:eq(0)>tr th:eq("+i+")", $curr).width( $(this).width());
});
//check to see if the width is set to auto, if not, we don't need to call the resizer function
if (options.width == "100%" || "auto"){
// call the resizer function whenever the table width has been adjusted
$(window).resize(function(){
resizer($this);
});
}
});
};
// private function to temporarily hide the header when the browser is resized
function resizer($this) {
// Need a relative reference to the "sticky header"
$curr = $this.prev();
$("thead:eq(0)>tr th", $this).each( function (i) {
$("thead:eq(0)>tr th:eq("+i+")", $curr).width( $(this).width());
});
};
})(jQuery);
And this is the corresponding html code:
Code:
<!-- DC DataGrid Start -->
<div id="tsort-tablewrapper" style="width: 100%;">
<div id="tsort-tableheader">
<div class="tsort-search">
<select id="tsort-columns" onchange="sorter.search('query')"> </select>
<input id="query" onkeyup="sorter.search('query')" type="text" />
</div>
<span class="tsort-details">
</span>
<div>
<span class="tsort-details">Records <span id="tsort-startrecord"></span>-<span id="tsort-endrecord"></span> of <span id="tsort-totalrecords"></span></span></div>
<span class="tsort-details"> </span>
<div>
<span class="tsort-details"><a href="javascript:sorter.reset()">reset</a></span>
</div>
</div>
<table border="0" cellspacing="10" class="tinytable" id="tsctablesort1"><thead>
<tr> <th><h3 id="col1">
ID</h3>
</th> <th><h3 id="col2">
MASTERPIECE</h3>
</th> <th><h3 id="col3">
ORIGIN</h3>
</th> <th><h3 id="col4">
COLOR</h3>
</th> </tr>
</thead> <tbody>
...
</tbody> </table>
<div id="tsort-tablefooter">
<div id="tsort-tablenav">
<div>
<img alt="START" height="16" onclick="sorter.move(-1,true)" src="dreamcodes/datagrid/images/first.gif" width="16" /> <img alt="PREVIOUS" height="16" onclick="sorter.move(-1)" src="dreamcodes/datagrid/images/previous.gif" width="16" /> <img alt="NEXT" height="16" onclick="sorter.move(1)" src="dreamcodes/datagrid/images/next.gif" width="16" /> <img alt="END" height="16" onclick="sorter.move(1,true)" src="dreamcodes/datagrid/images/last.gif" width="16" /> </div>
<div>
<select id="tsort-pagedropdown"> </select></div>
<div>
<a href="javascript:sorter.showall()">view all</a> </div>
</div>
<div id="tsort-tablelocation">
<div>
<select onchange="sorter.size(this.value)"> <option value="5">5</option> <option selected="selected" value="10">10</option> <option value="20">20</option> <option value="50">50</option> <option value="100">100</option> </select>
Entries Per Page </div>
<div class="tsort-page">
Page <span id="tsort-currentpage"></span> of <span id="tsort-totalpages"></span></div>
</div>
</div>
</div>
<!-- DC DataGrid Settings -->
<script type="text/javascript">
var sorter = new TINY.table.sorter('sorter','tsctablesort1',{
headclass:'head',
ascclass:'asc',
descclass:'desc',
evenclass:'tsort-evenrow',
oddclass:'tsort-oddrow',
evenselclass:'tsort-evenselected',
oddselclass:'tsort-oddselected',
paginate:true, // pagination (true,false)
size:10, // show 10 results per page
colddid:'tsort-columns',
currentid:'tsort-currentpage',
totalid:'tsort-totalpages',
startingrecid:'tsort-startrecord',
endingrecid:'tsort-endrecord',
totalrecid:'tsort-totalrecords',
hoverid:'tsort-selectedrow',
pageddid:'tsort-pagedropdown',
navid:'tsort-tablenav',
sortcolumn:1, // sort column 1
sortdir:1, // sort direction
columns:[{index:6, format:'%', decimals:1},{index:7, format:'$', decimals:0}], // classify for proper sorting
init:true // activate datagrid (true,false)
});
</script>
<!-- DC DataGrid End -->
<div class="tsc_clear">
</div>
<!-- line break/clear line -->
|
__label__pos
| 0.92119 |
FAQ: Am Fm Antenna Adapters?
How can I improve AM FM reception?
You can get better AM reception by using an external antenna such as a loop antenna, placing your radio near the window with its rear facing outside, or switching off/moving away any electrical appliances that have been proven to interfere with your radio’s reception.
Can I make my own FM antenna?
You can make an antenna for your FM radio with a 75 ohm wire, and attach it to the back of your tv. A good way to be able to listen to this without getting in trouble is by installing an antenna. To do this you will need a 75 ohm wire and some kind of plug that matches the end of the FM antenna on your tv.
Will a CB antenna work for AM FM radio?
You technically can use the CB antenna for your AM/FM radio, because your radio only needs something metallic and outside the vehicle.
How do I fix poor FM radio reception?
How to Fix Poor FM Radio Reception
1. Remove any obstacles you can.
2. Check and replace antenna connections.
3. Run a frequency scan.
4. Switch from stereo to mono.
5. Move your antenna: If you have an indoor antenna, place it near a window as high as possible to avoid interference from materials used in wall construction.
You might be interested: FAQ: Am Fm Transmittee?
Do FM antenna boosters work?
Antenna boosters can’t simply boost the signals broadcasted by the radio stations. It’s actually impossible. They work their wonders by increasing the gain that the antenna picks up.
What are the disadvantages of loop antenna?
Following are the disadvantages of Loop Antenna: ➨ Small loops have poor efficiency and hence are mainly used as receiving antenna at lower frequencies. ➨Small loop antennas have very low value of radiation resistance. This results into power loss as heat due to flow of current with high levels.
Does grounding antenna improve reception?
Does grounding the antenna improve the TV reception? Answer: It some rare case it may help the reception a little but usually it does not. Grounding will help reduce the damage and the chance for fire if the TV antenna were struck by lightning.
How can I make my FM antenna stronger?
Make an FM Dipole Antenna Design
1. Basic half wave dipole antenna.
2. Components required to make an FM dipole antenna.
3. FM dipole made from twin flex.
4. A tie wrap can be used to secure the centre point of the FM dipole antenna.
5. The wire end of the FM dipole is knotted.
How do I set up an FM antenna?
Insert the small end of the AC adapter into the FM antenna. Then insert the other end of the AC adapter into an outlet. Set the FM radio to a frequency near 98 MHz. Hold the FM antenna in your hands and move it around until you get the best reception so you know where to put it.
You might be interested: Question: How To Boost Fm Radio Signal Inside A Steel Building?
Can I use a digital TV antenna for FM radio?
Most digital television transmissions (including HDTV broadcasts) are in the UHF band (300 MHz to 3,000 MHz). But if it is a combined UHF/VHF antenna, with low-band VHF reception, then you should be able to receive the FM broadcasts.
Can I use a UHF antenna for FM radio?
Most UHF antennas will actually do a decent job of FM reception as well. UHF antennas don’t pick up the FM band by design, but they don’t block it either. It’s just that FM radio doesn’t require a very complex antenna for reception. In general, most antennas you’ll find for sale today are UHF/VHF antennas.
Is CB FM or AM?
CB mostly uses AM and amateur radio uses many different modes with FM being the most widely used mode above 30 MHz.
Are Firestik antennas good?
It’s our top recommendation if you’re looking for a quality, high performance CB antenna that will give you great range and last. The antenna has an adjustable tip, making it a snap to tune for optimal performance after installation. Trust us on this one: it’s hard to go wrong with the Firestik FS.
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|
__label__pos
| 0.999986 |
Ignore:
Timestamp:
Nov 10, 2003, 12:34:24 AM (17 years ago)
Author:
Sam Hocevar
Message:
• src/*.c: + Added missing #include "config.h". + Fixed a nasty bug in ee_rand() calls.
File:
1 edited
Legend:
Unmodified
Added
Removed
• ttyvaders/trunk/src/starfield.c
r90 r110
2121 */
2222
23#include "config.h"
24
2325#include <stdlib.h>
2426
3436 for(i = 0; i < STARS; i++)
3537 {
36 s[i].x = ee_rand(0, g->w);
37 s[i].y = ee_rand(0, g->h);
38 s[i].z = ee_rand(1, 4);
39 s[i].c = ee_rand(6, 8);
40 s[i].ch = ee_rand(0, 2) ? '.' : '\'';
38 s[i].x = ee_rand(0, g->w - 1);
39 s[i].y = ee_rand(0, g->h - 1);
40 s[i].z = ee_rand(1, 3);
41 s[i].c = ee_rand(6, 7);
42 s[i].ch = ee_rand(0, 1) ? '.' : '\'';
4143 }
4244
6769 if(s[i].x < 0)
6870 {
69 s[i].x = ee_rand(0, g->w);
71 s[i].x = ee_rand(0, g->w - 1);
7072 s[i].y = 0;
71 s[i].z = ee_rand(1, 3);
72 s[i].c = ee_rand(6, 8);
73 s[i].ch = ee_rand(0, 2) ? '.' : '\'';
73 s[i].z = ee_rand(1, 2);
74 s[i].c = ee_rand(6, 7);
75 s[i].ch = ee_rand(0, 1) ? '.' : '\'';
7476 }
7577 else if(s[i].y < g->h-1)
Note: See TracChangeset for help on using the changeset viewer.
|
__label__pos
| 0.982752 |
Tabata - November 29th, 2023
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Equipment:
- Weights
Tabata is a form of HIIT (High Intensity Interval Training) that challenges the participant to push themselves to their limits. It consists of 20 second "work" periods and 10 second "rest" periods repeated 8 times though for a total of 4 minutes. In Cat's Tabata class, we do 5 sets of Tabata with a minute rest between. Plus a 5 minute warmup for a total of a 30 minute workout.
Comments
|
__label__pos
| 0.503013 |
RT Journal Article SR Electronic T1 Multisystemic cellular tropism of SARS-CoV-2 in autopsies of COVID-19 patients JF medRxiv FD Cold Spring Harbor Laboratory Press SP 2021.06.03.21258241 DO 10.1101/2021.06.03.21258241 A1 Wong, Dickson W.L. A1 Klinkhammer, Barbara M. A1 Djudjaj, Sonja A1 Villwock, Sophia A1 Timm, M. Cherelle A1 Buhl, Eva M. A1 Wucherpfennig, Sophie A1 Cacchi, Claudio A1 Braunschweig, Till A1 Knüchel-Clarke, Ruth A1 Jonigk, Danny A1 Werlein, Christopher A1 Bülow, Roman D. A1 Dahl, Edgar A1 von Stillfried, Saskia A1 Boor, Peter YR 2021 UL http://medrxiv.org/content/early/2021/06/05/2021.06.03.21258241.abstract AB Background Multiorgan tropism of SARS-CoV-2 has previously been shown for several major organs.Methods We have comprehensively analyzed 25 different formalin-fixed paraffin-embedded (FFPE) tissues/organs from autopsies of fatal COVID-19 cases (n=8), using detailed histopathological assessment, detection of SARS-CoV-2 RNA using polymerase chain reaction and RNA in situ hybridization, viral protein using immunohistochemistry, and virus particles using transmission electron microscopy. Finally, we confirmed these findings in an independent external autopsy cohort (n=9).Findings SARS-CoV-2 RNA was mainly localized in epithelial cells, endothelial and mesenchymal cells across all organs. Next to lung, trachea, kidney, heart, or liver, viral RNA was also found in tonsils, salivary glands, oropharynx, thyroid, adrenal gland, testicles, prostate, ovaries, small bowel, lymph nodes, skin and skeletal muscle. Viral RNA was predominantly found in cells expressing ACE2, TMPRSS2, or both. The SARS-CoV-2 replicating RNA was also detected in these organs. Immunohistochemistry and electron microscopy were not suitable for reliable and specific SARS-CoV-2 detection in autopsies. The findings were validated using in situ hybridization on external COVID-19 autopsy samples. Finally, apart from the lung, correlation of virus detection and histopathological assessment did not reveal any specific alterations that could be attributed to SARS-CoV-2.Interpretation SARS-CoV-2 could be observed in virtually all organs, colocalizing with ACE2 and TMPRSS2 mainly in epithelial but also in mesenchymal and endothelial cells, and viral replication was found across all organ systems. Apart from the respiratory tract, no specific (histo-)morphologic alterations could be assigned to the SARS-CoV-2 infection.Evidence before this study SARS-CoV-2 has been shown to infect the respiratory tract and affect several other major organs. However, on a cellular level, the localization of SARS-CoV-2 and its targets ACE2 and TMPRSS2 have not been described comprehensively.Added value of this study We have analyzed tissue SARS-CoV-2 RNA using RT-PCR and visualized its localization together with ACE2 and TMPRSS2 using in situ hybridization (ISH) in 25 different autopsy tissues. SARS-CoV-2 sense and antisense RNA were detected in 16 tissues/organs, mainly in epithelial cells and, to a lesser extent, in endothelial or stromal cells. Detection of viral protein using immunohistochemistry or viral particles using transmission electron microscopy did not yield specific results. Interestingly, apart from the respiratory tract and specifically the lungs, we have not found a specific pathology that would be associated with extrapulmonary viral spread.Implications of all the available evidence We provide a recommendation on using these methods in autopsy diagnostics for SARS-CoV-2. Our data extend the current hypothesis of severe COVID-19 being multisystemic diseases. Our data also provide clear evidence of infection and replication of SARS-CoV-2 in the endothelial cell across all organs, extending the hypothesis on the (micro)vascular involvement in COVID-19.Competing Interest StatementThe authors have declared no competing interest.Funding StatementThis work was supported by the German Registry of COVID-19 Autopsies (www.DeRegCOVID.ukaachen.de), funded Federal Ministry of Health (ZMVI1-2520COR201), by the Federal Ministry of Education and Research within the framework of the network of university medicine (DEFEAT PANDEMICs, 01KX2021 and STOP-FSGS-01GM1901A), the German Research Foundation (DFG; SFB/TRR219 Project-ID 322900939, BO3755/13-1 Project-ID 454024652, and DJ100/1-1 432698239), the European Research Council (ERC) under the European Union's Horizon 2020 research and innovation program (grant agreement No 101001791), RWTH START-program (125/17 and 109/20) and the grant of the European Research Council (ERC); European Consolidator Grant, XHale to Danny Jonigk (ref. no.771883).Author DeclarationsI confirm all relevant ethical guidelines have been followed, and any necessary IRB and/or ethics committee approvals have been obtained.YesThe details of the IRB/oversight body that provided approval or exemption for the research described are given below:The study was given ethical approval by the ethics committee at the medical faculty of RWTH Aachen University (EK 304/20, EK 119/20, EK 092/20 and 9621_BO_K_2021).All necessary patient/participant consent has been obtained and the appropriate institutional forms have been archived.YesI understand that all clinical trials and any other prospective interventional studies must be registered with an ICMJE-approved registry, such as ClinicalTrials.gov. I confirm that any such study reported in the manuscript has been registered and the trial registration ID is provided (note: if posting a prospective study registered retrospectively, please provide a statement in the trial ID field explaining why the study was not registered in advance).YesI have followed all appropriate research reporting guidelines and uploaded the relevant EQUATOR Network research reporting checklist(s) and other pertinent material as supplementary files, if applicable.YesThe datasets generated and analyzed within the current study are not publicly available due to protecting the study participant privacy but are available from the corresponding author on reasonable request.
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__label__pos
| 0.601149 |
5 min read
How can technological innovation help the M&E sector deliver projects on time?
Tom Smith
Tom Smith
CAD Technician
Quite possibly the biggest change within construction over the past 60 years, has been the introduction of Computer-Aided Design (CAD). After the introduction of CAD, designing buildings became more efficient, accurate and faster. Building models that were once drawn by hand, are now designed on pieces of software within a matter of hours.
With technological innovations continuing to evolve within the built environment, can CAD help the M&E contractor achieve a quality M&E installation on time and within budget?
In this podcast we speak to two senior design CAD engineers, Skye Jennings and Tom Smith to answer just that.
We also look to provide answers to the following questions:
• What design software is typically used and how does it help an M&E contractor?
• Why do M&E designs go through so many iterations?
• Is it possible for BIM to be the solution?
• How would communication and early engagement help reduce design iterations?
• How can M&E contractors help the process?
• Are there any emerging technologies that could impact the design of service supports?
• Are there any unseen external factors that are commonly overlooked which may impact a full design?
Q1.) What design software is typically used and how does it help an M&E contractor?
As we are primarily working in 3D space, we draw our models in 3D CAD. We usually find that M&E contractors provide drawings in a variety of software – AutoCAD, Autodesk Inventor, Solidworks, Revitt, but often once an M&E project is over a certain size then a Navisworks model becomes the industry standard.
Revitt and Navisworks software have made the biggest impact on the quality and delivery time of an M&E installation. For the contractor on-site, it is best to work with Navisworks models because in most cases these models enable us to see everything in the surrounding environment to detect clashes. Clash checking is commonplace within an M&E installation, therefore a Navisworks model helps us to design and manufacture a support that can be installed in the correct place. Ultimately, accurate models reduce the likelihood of multiple design alterations shortening the time it takes for an M&E contractor to receive our support frames.
Q2.) Why do M&E designs go through so many iterations?
This is due to the nature of an M&E project. Typically, there are many different contractors on-site with different installation deadlines, on top of this, commercial factors such as price increases and supply chain volatility can mean different components are purchased, resulting in design changes.
To give an example, a certain brand of AHU originally specified for a project may not be available due to commercial issues, this would mean a complete change in the design with a change in hole patterns, fixings, sizes and weights.
Q3.) Will BIM help to reduce the time between design and delivery?
The UK government believes that one way to improve the efficiency, cost and safety of construction projects within the UK is by implementing a BIM framework (Construction Playbook 2020). With many within the industry aspiring to work within a BIM framework, it is unlikely to have the necessary positive impact without better engagement between all contractors on-site and their respective supply chains.
Q4.) How would communication and early engagement help reduce design iterations?
Early and continued engagement with the contractor is critical. Installers are often four or even five tiers down from the main contractor and they may not even have a Navisworks drawing to work too, this will mean that they won’t know how other services would interfere with their installation. If there is free-flowing communication between the principal contractor, the architect, main contractor, subcontractors (installers) and the M&E support services supplier then any likely changes on-site can be communicated earlier on in the process so design iterations are not required further down the line.
Project timeline
Q5.) How can M&E contractors help the process?
M&E contractors can help us reduce the time to create a design, by firstly providing us with a suitable model and secondly by communicating with us any changes on-site that would impact that design.
In a worst-case scenario we receive a ‘fag packet’ drawing. A ‘fag packet’ drawing is not fit for purpose, as it isn’t a 3D model, it is unlikely to have accurate dimensions or critical information such as load paths which would impact the bracket, we are designing supports for. We would always prefer a Navisworks, Revit, DWG or step file drawing.
Our preference would be a Navisworks model as they take into consideration the whole environment enabling us to detect clashes and meet any necessary regulations. This will allow an M&E installation to be completed in a much shorter time.
Navisworks model
Navisworks model example
Q6.) Are there any emerging technological innovations that could impact the design of service supports?
There are several exciting types of technology utilised within the industry. For instance, highly accurate 3D scanning using augmented and virtual reality integrating with BIM can be invaluable. For example, a few sites use AR with Revitt and are trialing the software with Microsoft HoloLens.
This technology enables the user to visualise, within a couple of millimeters accuracy, the location of services before they are installed, resulting in more accurate designs and installations. The main barrier to this technology is cost, as installers on-site are unlikely to have the finances available to purchase the equipment and therefore won’t be able to benefit from this technology.
In one project we had a Revitt Navisworks model but it wasn’t accurate. The site had concrete dishes but had been set differently on site to what the model had shown. This was critical as we were fixing a load of supports onto concrete, after we managed to receive a full 3D scan of the site, and imported that into Inventor, we could accurately design a support to the mm, when you need this level of accuracy, the technology is invaluable.
Q7.) Are there any unseen external factors that are commonly overlooked which may impact a full design?
As the M&E industry becomes ever more litigious, compliant installations matter. To be compliant, evidence is required to show that any support supplied to site is fit for purpose. For an M&E contractor having that evidence is important, if, factors on-site change, for example, with a last-minute component change then a new design must be drawn. With the increasingly tight deadlines for project delivery, M&E contractors will be looking to try to limit any changes that could result in a design iteration.
Parasitic loading
Parasitic loading
Another common issue within the M&E industry is parasitic loading, this is where an installer comes in after a bracket has been installed and fixes (piggybacks) more supports of the bracket. This is a problem because once a bracket has been signed off to safely support the load-bearing capacity of a particular service load and then more services are added to the bracket, this immediately compromises the load bearing capacity. In this situation, a new bracket must be made, with a separate load calculation to take into consideration the extra load.
Lack of early engagement
In many cases, the complexity and limited communication with the whole supply chain often leads to design iterations. With complex M&E installations, the earlier the supply chain engages with a project, the more likely supports can be designed and manufactured fit for purpose.
What are the top requirements from a design perspective?
• Early and continued engagement so that design changes can be made where necessary.
• An accurate model of the building showing the whole environment to detect clashes.
• Having up-to-date and accurate information to hand.
Conclusion
There will always be external influences, dictated by market conditions, that will impact the delivery of a project. However, with early and continued engagement between different stakeholders within the supply chain, and appropriate design models to work from, accurate designs can be produced, enabling M&E contractors to deliver projects on time.
Topics: DfMA
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__label__pos
| 0.904061 |
Education
Understanding Bipolar Disorder
Bipolar Disorder
Bipolar Disorder
TNews – Bipolar disorder is a complex mental health condition characterized by extreme mood swings that can significantly impact an individual’s life. These mood swings fluctuate between two distinct poles: manic episodes and depressive episodes.
Manic Episodes:
During manic episodes, individuals experience heightened emotions, often to an extreme degree. They may feel an overwhelming sense of joy or euphoria that is far beyond what is considered normal. This elevated mood is accompanied by a surge in energy levels, making individuals feel highly motivated and confident. However, these episodes also bring with them a range of challenging symptoms, such as difficulty sleeping or insomnia, an inflated sense of self-importance, increased activity levels, rapid and constant speech, and engaging in impulsive and risky behaviors, like excessive spending.
Depressive Episodes:
Conversely, depressive episodes are marked by intense feelings of sadness, hopelessness, and worthlessness. Individuals with bipolar disorder often lose interest in activities they once enjoyed and struggle to find pleasure in life. Changes in appetite, either leading to significant weight loss or gain, are common during these episodes. Sleep disturbances, including insomnia, further exacerbate the situation, leaving individuals fatigued and drained of energy. Concentration and decision-making become arduous tasks, and, in severe cases, thoughts of death or suicide may emerge.
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__label__pos
| 0.962306 |
FFmpeg coverage
Directory: ../../../ffmpeg/
File: src/libavcodec/x86/me_cmp_init.c
Date: 2024-04-19 17:50:32
Exec Total Coverage
Lines: 66 69 95.7%
Functions: 7 7 100.0%
Branches: 19 24 79.2%
Line Branch Exec Source
1 /*
2 * SIMD-optimized motion estimation
3 * Copyright (c) 2000, 2001 Fabrice Bellard
4 * Copyright (c) 2002-2004 Michael Niedermayer <[email protected]>
5 *
6 * MMX optimization by Nick Kurshev <[email protected]>
7 *
8 * This file is part of FFmpeg.
9 *
10 * FFmpeg is free software; you can redistribute it and/or
11 * modify it under the terms of the GNU Lesser General Public
12 * License as published by the Free Software Foundation; either
13 * version 2.1 of the License, or (at your option) any later version.
14 *
15 * FFmpeg is distributed in the hope that it will be useful,
16 * but WITHOUT ANY WARRANTY; without even the implied warranty of
17 * MERCHANTABILITY or FITNESS FOR A PARTICULAR PURPOSE. See the GNU
18 * Lesser General Public License for more details.
19 *
20 * You should have received a copy of the GNU Lesser General Public
21 * License along with FFmpeg; if not, write to the Free Software
22 * Foundation, Inc., 51 Franklin Street, Fifth Floor, Boston, MA 02110-1301 USA
23 */
24
25 #include "libavutil/attributes.h"
26 #include "libavutil/cpu.h"
27 #include "libavutil/mem_internal.h"
28 #include "libavutil/x86/asm.h"
29 #include "libavutil/x86/cpu.h"
30 #include "libavcodec/me_cmp.h"
31 #include "libavcodec/mpegvideo.h"
32
33 int ff_sum_abs_dctelem_sse2(const int16_t *block);
34 int ff_sum_abs_dctelem_ssse3(const int16_t *block);
35 int ff_sse8_mmx(MpegEncContext *v, const uint8_t *pix1, const uint8_t *pix2,
36 ptrdiff_t stride, int h);
37 int ff_sse16_mmx(MpegEncContext *v, const uint8_t *pix1, const uint8_t *pix2,
38 ptrdiff_t stride, int h);
39 int ff_sse16_sse2(MpegEncContext *v, const uint8_t *pix1, const uint8_t *pix2,
40 ptrdiff_t stride, int h);
41 int ff_hf_noise8_mmx(const uint8_t *pix1, ptrdiff_t stride, int h);
42 int ff_hf_noise16_mmx(const uint8_t *pix1, ptrdiff_t stride, int h);
43 int ff_sad8_mmxext(MpegEncContext *v, const uint8_t *pix1, const uint8_t *pix2,
44 ptrdiff_t stride, int h);
45 int ff_sad16_mmxext(MpegEncContext *v, const uint8_t *pix1, const uint8_t *pix2,
46 ptrdiff_t stride, int h);
47 int ff_sad16_sse2(MpegEncContext *v, const uint8_t *pix1, const uint8_t *pix2,
48 ptrdiff_t stride, int h);
49 int ff_sad8_x2_mmxext(MpegEncContext *v, const uint8_t *pix1, const uint8_t *pix2,
50 ptrdiff_t stride, int h);
51 int ff_sad16_x2_mmxext(MpegEncContext *v, const uint8_t *pix1, const uint8_t *pix2,
52 ptrdiff_t stride, int h);
53 int ff_sad16_x2_sse2(MpegEncContext *v, const uint8_t *pix1, const uint8_t *pix2,
54 ptrdiff_t stride, int h);
55 int ff_sad8_y2_mmxext(MpegEncContext *v, const uint8_t *pix1, const uint8_t *pix2,
56 ptrdiff_t stride, int h);
57 int ff_sad16_y2_mmxext(MpegEncContext *v, const uint8_t *pix1, const uint8_t *pix2,
58 ptrdiff_t stride, int h);
59 int ff_sad16_y2_sse2(MpegEncContext *v, const uint8_t *pix1, const uint8_t *pix2,
60 ptrdiff_t stride, int h);
61 int ff_sad8_approx_xy2_mmxext(MpegEncContext *v, const uint8_t *pix1, const uint8_t *pix2,
62 ptrdiff_t stride, int h);
63 int ff_sad16_approx_xy2_mmxext(MpegEncContext *v, const uint8_t *pix1, const uint8_t *pix2,
64 ptrdiff_t stride, int h);
65 int ff_sad16_approx_xy2_sse2(MpegEncContext *v, const uint8_t *pix1, const uint8_t *pix2,
66 ptrdiff_t stride, int h);
67 int ff_vsad_intra8_mmxext(MpegEncContext *v, const uint8_t *pix1, const uint8_t *pix2,
68 ptrdiff_t stride, int h);
69 int ff_vsad_intra16_mmxext(MpegEncContext *v, const uint8_t *pix1, const uint8_t *pix2,
70 ptrdiff_t stride, int h);
71 int ff_vsad_intra16_sse2(MpegEncContext *v, const uint8_t *pix1, const uint8_t *pix2,
72 ptrdiff_t stride, int h);
73 int ff_vsad8_approx_mmxext(MpegEncContext *v, const uint8_t *pix1, const uint8_t *pix2,
74 ptrdiff_t stride, int h);
75 int ff_vsad16_approx_mmxext(MpegEncContext *v, const uint8_t *pix1, const uint8_t *pix2,
76 ptrdiff_t stride, int h);
77 int ff_vsad16_approx_sse2(MpegEncContext *v, const uint8_t *pix1, const uint8_t *pix2,
78 ptrdiff_t stride, int h);
79
80 #define hadamard_func(cpu) \
81 int ff_hadamard8_diff_ ## cpu(MpegEncContext *s, const uint8_t *src1, \
82 const uint8_t *src2, ptrdiff_t stride, int h); \
83 int ff_hadamard8_diff16_ ## cpu(MpegEncContext *s, const uint8_t *src1, \
84 const uint8_t *src2, ptrdiff_t stride, int h);
85
86 hadamard_func(mmxext)
87 hadamard_func(sse2)
88 hadamard_func(ssse3)
89
90 #if HAVE_X86ASM
91 16 static int nsse16_mmx(MpegEncContext *c, const uint8_t *pix1, const uint8_t *pix2,
92 ptrdiff_t stride, int h)
93 {
94 int score1, score2;
95
96
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16 if (c)
97 score1 = c->mecc.sse[0](c, pix1, pix2, stride, h);
98 else
99 16 score1 = ff_sse16_mmx(c, pix1, pix2, stride, h);
100 16 score2 = ff_hf_noise16_mmx(pix1, stride, h) + ff_hf_noise8_mmx(pix1+8, stride, h)
101 16 - ff_hf_noise16_mmx(pix2, stride, h) - ff_hf_noise8_mmx(pix2+8, stride, h);
102
103
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16 if (c)
104 return score1 + FFABS(score2) * c->avctx->nsse_weight;
105 else
106 16 return score1 + FFABS(score2) * 8;
107 }
108
109 16 static int nsse8_mmx(MpegEncContext *c, const uint8_t *pix1, const uint8_t *pix2,
110 ptrdiff_t stride, int h)
111 {
112 16 int score1 = ff_sse8_mmx(c, pix1, pix2, stride, h);
113 16 int score2 = ff_hf_noise8_mmx(pix1, stride, h) -
114 16 ff_hf_noise8_mmx(pix2, stride, h);
115
116
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16 if (c)
117 return score1 + FFABS(score2) * c->avctx->nsse_weight;
118 else
119 16 return score1 + FFABS(score2) * 8;
120 }
121
122 #endif /* HAVE_X86ASM */
123
124 #if HAVE_INLINE_ASM
125
126 DECLARE_ASM_CONST(8, uint64_t, round_tab)[3] = {
127 0x0000000000000000ULL,
128 0x0001000100010001ULL,
129 0x0002000200020002ULL,
130 };
131
132 71328 static inline void sad8_4_mmx(const uint8_t *blk1, const uint8_t *blk2,
133 ptrdiff_t stride, int h)
134 {
135 71328 x86_reg len = -stride * h;
136 71328 __asm__ volatile (
137 "movq (%1, %%"FF_REG_a"), %%mm0\n\t"
138 "movq 1(%1, %%"FF_REG_a"), %%mm2\n\t"
139 "movq %%mm0, %%mm1 \n\t"
140 "movq %%mm2, %%mm3 \n\t"
141 "punpcklbw %%mm7, %%mm0 \n\t"
142 "punpckhbw %%mm7, %%mm1 \n\t"
143 "punpcklbw %%mm7, %%mm2 \n\t"
144 "punpckhbw %%mm7, %%mm3 \n\t"
145 "paddw %%mm2, %%mm0 \n\t"
146 "paddw %%mm3, %%mm1 \n\t"
147 ".p2align 4 \n\t"
148 "1: \n\t"
149 "movq (%2, %%"FF_REG_a"), %%mm2\n\t"
150 "movq 1(%2, %%"FF_REG_a"), %%mm4\n\t"
151 "movq %%mm2, %%mm3 \n\t"
152 "movq %%mm4, %%mm5 \n\t"
153 "punpcklbw %%mm7, %%mm2 \n\t"
154 "punpckhbw %%mm7, %%mm3 \n\t"
155 "punpcklbw %%mm7, %%mm4 \n\t"
156 "punpckhbw %%mm7, %%mm5 \n\t"
157 "paddw %%mm4, %%mm2 \n\t"
158 "paddw %%mm5, %%mm3 \n\t"
159 "movq %5, %%mm5 \n\t"
160 "paddw %%mm2, %%mm0 \n\t"
161 "paddw %%mm3, %%mm1 \n\t"
162 "paddw %%mm5, %%mm0 \n\t"
163 "paddw %%mm5, %%mm1 \n\t"
164 "movq (%3, %%"FF_REG_a"), %%mm4 \n\t"
165 "movq (%3, %%"FF_REG_a"), %%mm5 \n\t"
166 "psrlw $2, %%mm0 \n\t"
167 "psrlw $2, %%mm1 \n\t"
168 "packuswb %%mm1, %%mm0 \n\t"
169 "psubusb %%mm0, %%mm4 \n\t"
170 "psubusb %%mm5, %%mm0 \n\t"
171 "por %%mm4, %%mm0 \n\t"
172 "movq %%mm0, %%mm4 \n\t"
173 "punpcklbw %%mm7, %%mm0 \n\t"
174 "punpckhbw %%mm7, %%mm4 \n\t"
175 "paddw %%mm0, %%mm6 \n\t"
176 "paddw %%mm4, %%mm6 \n\t"
177 "movq %%mm2, %%mm0 \n\t"
178 "movq %%mm3, %%mm1 \n\t"
179 "add %4, %%"FF_REG_a" \n\t"
180 " js 1b \n\t"
181 : "+a" (len)
182 71328 : "r" (blk1 - len), "r" (blk1 - len + stride), "r" (blk2 - len),
183 "r" (stride), "m" (round_tab[2]));
184 71328 }
185
186 35672 static inline int sum_mmx(void)
187 {
188 int ret;
189 35672 __asm__ volatile (
190 "movq %%mm6, %%mm0 \n\t"
191 "psrlq $32, %%mm6 \n\t"
192 "paddw %%mm0, %%mm6 \n\t"
193 "movq %%mm6, %%mm0 \n\t"
194 "psrlq $16, %%mm6 \n\t"
195 "paddw %%mm0, %%mm6 \n\t"
196 "movd %%mm6, %0 \n\t"
197 : "=r" (ret));
198 35672 return ret & 0xFFFF;
199 }
200
201 #define PIX_SADXY(suf) \
202 static int sad8_xy2_ ## suf(MpegEncContext *v, const uint8_t *blk2, \
203 const uint8_t *blk1, ptrdiff_t stride, int h) \
204 { \
205 __asm__ volatile ( \
206 "pxor %%mm7, %%mm7 \n\t" \
207 "pxor %%mm6, %%mm6 \n\t" \
208 ::); \
209 \
210 sad8_4_ ## suf(blk1, blk2, stride, h); \
211 \
212 return sum_ ## suf(); \
213 } \
214 \
215 static int sad16_xy2_ ## suf(MpegEncContext *v, const uint8_t *blk2, \
216 const uint8_t *blk1, ptrdiff_t stride, int h) \
217 { \
218 __asm__ volatile ( \
219 "pxor %%mm7, %%mm7 \n\t" \
220 "pxor %%mm6, %%mm6 \n\t" \
221 ::); \
222 \
223 sad8_4_ ## suf(blk1, blk2, stride, h); \
224 sad8_4_ ## suf(blk1 + 8, blk2 + 8, stride, h); \
225 \
226 return sum_ ## suf(); \
227 } \
228
229 71344 PIX_SADXY(mmx)
230
231 #endif /* HAVE_INLINE_ASM */
232
233 1072 av_cold void ff_me_cmp_init_x86(MECmpContext *c, AVCodecContext *avctx)
234 {
235 1072 int cpu_flags = av_get_cpu_flags();
236
237 #if HAVE_INLINE_ASM
238
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1072 if (INLINE_MMX(cpu_flags)) {
239 68 c->pix_abs[0][3] = sad16_xy2_mmx;
240 68 c->pix_abs[1][3] = sad8_xy2_mmx;
241 }
242
243 #endif /* HAVE_INLINE_ASM */
244
245
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1072 if (EXTERNAL_MMX(cpu_flags)) {
246 68 c->sse[1] = ff_sse8_mmx;
247 #if HAVE_X86ASM
248 68 c->nsse[0] = nsse16_mmx;
249 68 c->nsse[1] = nsse8_mmx;
250 #endif
251 }
252
253
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1072 if (EXTERNAL_MMXEXT(cpu_flags)) {
254 #if !HAVE_ALIGNED_STACK
255 c->hadamard8_diff[0] = ff_hadamard8_diff16_mmxext;
256 c->hadamard8_diff[1] = ff_hadamard8_diff_mmxext;
257 #endif
258
259 67 c->sad[0] = ff_sad16_mmxext;
260 67 c->sad[1] = ff_sad8_mmxext;
261
262 67 c->pix_abs[0][0] = ff_sad16_mmxext;
263 67 c->pix_abs[0][1] = ff_sad16_x2_mmxext;
264 67 c->pix_abs[0][2] = ff_sad16_y2_mmxext;
265 67 c->pix_abs[1][0] = ff_sad8_mmxext;
266 67 c->pix_abs[1][1] = ff_sad8_x2_mmxext;
267 67 c->pix_abs[1][2] = ff_sad8_y2_mmxext;
268
269 67 c->vsad[4] = ff_vsad_intra16_mmxext;
270 67 c->vsad[5] = ff_vsad_intra8_mmxext;
271
272
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67 if (!(avctx->flags & AV_CODEC_FLAG_BITEXACT)) {
273 54 c->pix_abs[0][3] = ff_sad16_approx_xy2_mmxext;
274 54 c->pix_abs[1][3] = ff_sad8_approx_xy2_mmxext;
275
276 54 c->vsad[0] = ff_vsad16_approx_mmxext;
277 54 c->vsad[1] = ff_vsad8_approx_mmxext;
278 }
279 }
280
281
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1072 if (EXTERNAL_SSE2(cpu_flags)) {
282 65 c->sse[0] = ff_sse16_sse2;
283 65 c->sum_abs_dctelem = ff_sum_abs_dctelem_sse2;
284
285 #if HAVE_ALIGNED_STACK
286 65 c->hadamard8_diff[0] = ff_hadamard8_diff16_sse2;
287 65 c->hadamard8_diff[1] = ff_hadamard8_diff_sse2;
288 #endif
289
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65 if (!(cpu_flags & AV_CPU_FLAG_SSE2SLOW) && avctx->codec_id != AV_CODEC_ID_SNOW) {
290 65 c->sad[0] = ff_sad16_sse2;
291 65 c->pix_abs[0][0] = ff_sad16_sse2;
292 65 c->pix_abs[0][1] = ff_sad16_x2_sse2;
293 65 c->pix_abs[0][2] = ff_sad16_y2_sse2;
294
295 65 c->vsad[4] = ff_vsad_intra16_sse2;
296
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65 if (!(avctx->flags & AV_CODEC_FLAG_BITEXACT)) {
297 54 c->pix_abs[0][3] = ff_sad16_approx_xy2_sse2;
298 54 c->vsad[0] = ff_vsad16_approx_sse2;
299 }
300 }
301 }
302
303
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1072 if (EXTERNAL_SSSE3(cpu_flags)) {
304 63 c->sum_abs_dctelem = ff_sum_abs_dctelem_ssse3;
305 #if HAVE_ALIGNED_STACK
306 63 c->hadamard8_diff[0] = ff_hadamard8_diff16_ssse3;
307 63 c->hadamard8_diff[1] = ff_hadamard8_diff_ssse3;
308 #endif
309 }
310 1072 }
311
|
__label__pos
| 0.977763 |
4
$\begingroup$
I have a bunch of questions concerning the use of the copula package in R. My overall aim is to generate synthetic values using copulas. I am analyzing a hydrological data: annual peak discharge [m³/s] and corresponding volume [m³].
I managed to apply tests on serial independence and dependence. Furthermore I identified and excluded ties and created pseudo-observations (transformation of copula values between [0,1]). Since I don’t know which copula is the best, I fitted the copula parameter first:
fg <- fitCopula(copula=gumbelCopula(), data=u) # u is my data
I will do this for all available copulas in R. Afterwards I test the goodness of fit with the following function:
gofCopula(copula=gumbelCopula(fg@estimate), x=data[,2:3], N=1000, method="Sn",
estim.method="mpl", simulation="mult")
Using the “best” copula, I then want to create synthetic values. I found a function to create random samples, but I am not sure, if it does what I need.
random_samples <- rCopula(copula=gumbelCopula(fg@estimate), n=10000)
It seems to me that this function creates only random values, but is the dependency structure of my data set considered? There is also another function in the copula package mvdc, for the construction of multivariate distributions from copulas. What is actually the difference of mvdc and rCopula, both are generating synthetic values, aren’t they?
One last question is: Once I am able to generate my synthetic values, how can I transform them back to their real units? From reading through the documentation I understood tat I have to multiply the values for (u,v) with the inverse of their particular cdf, is this true?
One question is not answered yet, I want to render it more precisely: namely the function . According to the copula manual p. 107, it is used to "construct multivariate distributions from copulas"
For the function a copula family, as well as the distributions of the margins have to be specified, for instance:
mv.NE <- mvdc(copula=gumbelCopula(fg@estimate),margins=c("norm","norm"),paramMargins=list(list(mean=0, sd = 1),list(mean=0, sd = 1)))
(here is chose a gumbelCopula and estimated the parameter with . I assume that my marginals are both distributed "normal".)
What for do I need this function? I am slightly confused because of:
• when I create random values from my copula using I eventually yield "distributed" values when I pass the values of u and v respectively over to their particular distribution function (in this case both are distributed normal)
• so why is there a second option to create multivariate distributions from a copula.
I just don't get the difference...
$\endgroup$
3
• $\begingroup$ Correct, after you sample from the copula, you transform the margins individually, in the way you suggest. $\endgroup$
– Glen_b
Mar 20 '14 at 23:42
• $\begingroup$ Yes, the random samples that are generated are not independent; they contain the dependence structure that results from the copula (and parameter estimate(s)) you supply. You can see they're dependent simply by plotting the result. $\endgroup$
– Glen_b
Mar 20 '14 at 23:59
• $\begingroup$ Just to be clear, you do not have to multiply the values of (u,v) with their inverse marignal CDFs, but you have to pass the values (u,v) to their inverse marginal CDFs. $\endgroup$
– Ben
Mar 21 '14 at 10:35
2
$\begingroup$
Thanks for your answers. Indeed, when I plot the different Copulas, it seems (based on a visual comparison) that the dependence struture is depicted in the copula.
Comparison of different results
Then my next step will be to compare the different copula families in terms of looking at their goodness of fit.
@Ben: Ok, so if the distribution of one of my marginals is the generalized extreme value distribution, my u or v will be passed to the "x" of this function, since the other parameters of this distribution are fitted before. forumla_gev
$\endgroup$
7
• 1
$\begingroup$ You might want to interactively explore different bivariate copula families that are available in R with the shiny app copulatheque. $\endgroup$
– Ben
Mar 21 '14 at 10:40
• $\begingroup$ You need to use the inverse cumulative distribution function (CDF) also known as quantile function and commonly abbreviated in R as "q" (e.q. qnorm, qgev, ...). This should be clear from Sklar's Theorem. Please pay careful attention to the theory of copulas, as it is very easy to mess around with them. $\endgroup$
– Ben
Mar 24 '14 at 9:48
• $\begingroup$ Hi Ben, sorry for the delay. Thanks for that hint! I am very glad that the GEV is already coded in R. Telling from the destribution functions on link I thought it wasn't. I will use this one. $\endgroup$
– Jochen
Mar 26 '14 at 12:15
• $\begingroup$ You will need an extra package: "evd" $\endgroup$
– Ben
Mar 27 '14 at 9:26
• $\begingroup$ @Ben thanks for linking to that app! It's incredibly helpful to visualize the copulas this way! $\endgroup$ Aug 4 '15 at 17:12
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How does every human body smell? Can this change?
How does every human body smell? Can this change?
When we meet someone, sometimes we even feel the special smell coming from his body. It is said that every person’s body odor is unique and special. For this reason, each person can be identified by smell. Based on this, the police search for sniffer dogs. It is true that every body odor is different. […]
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Tutorial 11: Random Numbers
Random Numbers...
HOME >> Tutorials >> Tutorial 11: Random Numbers
Introduction:
In this tutorial we will learn our last command, randInt(. This really has nothing to do with the rest of the commands that I have taught you. In fact it isn't even a PRGM command, but a MATH command. I just had to insert this last tutorial because you need the command for the game in the next tutorial. The command is fairly easy to grasp, because all it does is pick a random number (an integer) between the range that you set.
New Commands:
randInt( Generates and displays a random integer within a range
randInt(lower,upper[,numtrials])
The Codes:
There are going to be two separate codes: one shows you the basic way of using randInt in a program and the second is a fun little program that I like to put on people's calculator. However, before even getting to the codes, I want you to really grasp how randInt works. Before entering the program editor, I want you to go to the home screen and enter the randInt command. I want your calculator to display a random number between 1 and 50, so it should look like this randInt(1,50). I don't want any trials. Okay, press ENTER and you should see a number that's between 1 and 50. Now press ENTER again, and a new number between 1 and 50 should appear. If it so happened that the same number appeared, then the 1 in 50 chance happened and the randInt command pick the same number again. Just press enter again and there should be a new number. Now keep on pressing Í until you see the same number twice. I just want you to really see how the randInt command works. Now on to the codes:
: For AShell, SOS, and TI-Explorer
Lbl 1
0STOREB Stores 0 into X
Repeat A=19 Repeat everything until the random number is 19
ClrHome Clears the home screen
B+1STOREB Add one to X every time it repeats
randInt(1,20)STOREA Picks a random integer between 1 and 20
Output(4,8,A) Displays that random number on line 4
For(X,1,100) This construction is used so that we will be able
End to see the random number for a short while
End End Repeat Loop
Output(8,8,B) Displays the number of times it took to get to 19
An example of a final output:
This program just keeps on displaying random numbers between 1 and 20 until the random number is 19. If it so happens (it happened to me when I was testing the program) that the first random number is 19, just run the program again. User variable B in this program is the number of times it took for the random number to be 19. So if the number at the bottom (B) was 23, that means that 19 appeared on the 23rd try. You better study that because you'll need it in the next tutorial. The next program is a fun little program:
: For AShell, SOS, and TI-Explorer
Lbl 1
ClrHome Clears the home screen
randInt(1,8)STOREA Picks a random integer between 1 and 8
randInt(1,12)STOREB Picks a random integer between 1 and 12
Output(A,B,"HELLO") Displays "HELLO" in a random spot
For(X,1,100) This construction is used so that we will be able
End to see the random number for a short while
Goto 1 Go to Label 1
A possible output:
As you have probably noticed, the program doesn't stop. This is an endless loop (I did this on purpose) and it only stops unless you press ON and then ENTER or press 2ND and ON [OFF]. I was thinking of putting in a For statement that would let it display "HELLO" in a random spot for a certain number of times, but I decided to do it this way just so you would learn how to break a program. However, when I make this program I use a command called getKey, so it'll stop displaying "HELLO" if I either: press ON and then ENTER, press 2ND and ON [OFF], or I press a particular key. To take a look at how this is done, you can go to Tutorial 21.
Version 2.0 Update:
After the first output of BASIC 83 Guru, I thought of another way to use randInt. This way of using randInt is if you want more than one choice to happen randomly after a certain situation occurs. randInt is accompanied by Goto & Lbl and If statements in this code. What if you wanted to make a "Magic 8 Ball" game, which had four different replies. You would want those replies to be randomly chosen, right? Well this is how you'd make such a game:
The Code:
: For AShell, SOS, and TI-Explorer
Lbl A
ClrHome Clears the home screen
Output(4,3,"MAGIC 8 BALL") Displays text at row 4 and column 3
For(X,1,500) This line and the next line are used to delay
End the program for a couple of seconds
Lbl D
ClrHome Clears the home screen
Output(3,2,"THINK OF YOUR") Displays "THINK OF YOUR"
Output(4,3,"QUESTION AND") Outputs text at row 3 and column 3
Output(5,1,"PRESS ENTER FOR") Displays "PRESS ENTER FOR"
Output(6,4,"AN ANSWER.") Outputs text at row 6 and column 4
Pause Suspends program and waits for user to press ENTER
ClrHome Clears the home screen
Output(4,5,"THINKING") Displays "THINKING" on the screen for
For(X,1,500) screen for a few seconds
End
randInt(1,4)STOREC Picks a random integer between 1 and 4
ClrHome Clears the home screen
If C=1:Goto 1 If C equals 1, go to Label 1
If C=2:Goto 2 If C equals 2, go to Label 2
If C=3:Goto 3 If C equals 3, go to Label 3
If C=4:Goto 4 If C equals 4, go to Label 4
Lbl 1
Output(4,2,"DEFINITELY YES") Displays "DEFINITELY YES"
Goto B Goto Label B
Lbl 2
Output(4,2,"DEFINITELY NOT") Displays "DEFINITELY NOT"
Goto B Goto Label B
Lbl 3
Output(4,5,"MAYBE SO") Displays "MAYBE SO"
Goto B Goto Label B
Lbl 4
Output(4,1,"TRY AGAIN LATER") Displays "TRY AGAIN LATER"
Goto B Goto Label B
Lbl B
Pause Suspends program and waits for user to press ENTER
Menu("AGAIN?","YES",D,"NO",C) Ask if you want to play again
Lbl C
Possible output:
then
AND
and
Further Explanation:
If you've run the program you can see that the opening screen comes up displaying "MAGIC 8 BALL" and it waits a few seconds and then asks you to think of the question that you want to ask the "Magic 8 Ball." When you press ENTER it displays "THINKING" for a while and then gives you the answer. I added that "THINKING" part, just so it gives the effect that the game is actually thinking. If I didn't put that part there, it would almost immediately give me the answer, which isn't very realistic. All of the answers have a 1 in 4 (equal) chance of appearing. Now what happens if you wanted "MAYBE SO" to have a better chance of occurring than all the other answers? What if you wanted it to be weighted? This is really easy. All you have to do is change two lines. Change the line that reads randInt(1,4)STOREC to randInt(1,5)STOREC and change the line that reads If C=3:Goto 3 to If C=3 or C=5:Goto 3. This makes the chance that "MAYBE SO" is displayed 2 out of 5 and the others 1 out of 5. You may be wondering what the command or does. Will it's pretty simple actually. If either one of the statements are true, it'll go to label 3. You can find or by pressing 2ND, MATH, and then RIGHT. Scroll down to 2:or. This construction of randInt is used a lot in "text only" games or RPGs (role playing games). It makes the games better because different outcomes can occur randomly. If a certain outcome was fixed, the program would be boring after the first couple of times playing it.
Conclusion:
Well, this is the last command that you'll learn before our big game. If you want to quickly review all the lessons up to this point, just click on the "Tutorials" link in either navigation bar. If you feel confident that you completely understand everything that I have taught you, then go on to The Big Game.
If you do not understand a particular part in this lesson, have suggestions, or find any problems please contact me.
Previous Tutorial Next Tutorial
Tutorial 10 The Big Game
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实验室实验台水电如何进行安装？
来源：本站 发布时间：2022-02-24
1、实验室家具、通风柜、水电工程等安装条件
1.1 上下水：建筑施工单位应预先将上下水预留在我们指定的室内位置，上水管应带有阀门，以便我们进行安装。
1.2 强电、弱电：每个实验室都应建筑施工单位配备配电箱，应有足够的电流范围以及接口位置供我们引线，配电箱内应配备空开及漏电保护。每个实验室应配备仪器所需要的所有Lan线接口，由我方统一引入实验台。
1.3 墙体、地面、天花：在我方施工前，建筑施工单位应完成地面、墙壁与买方的交接。地面平整度误差小于3mm 。
2、实验台安装
2.1 实验桌落地型底柜:
a) 将各单元底柜按图标位置直立摆放并调整水平，将各单元底柜以同样的水平位置连结，如为独立之单元底柜亦应予以适当的方式固定于实验室地板上。
b) 安装配件，并注意各类公用系统管线不可有任何的缠绕或纠结。
c) 如有需要，各单元将以同一水平做紧密结合，邻接的抽屉及门片均在同一中线。
d) 调整各底柜之柜门抽屉以可顺利开关，不会有任何阻碍或丧失正常功能。
2.2 实验桌台面:
a) 现场施工： 当台面板接缝须于施工现场黏接时，其工序应与工厂内黏接方式相同，必须使用生产厂商规定之专用黏合剂，以及生产厂商所规定之夹具及施工方式做连接，台面板的大部裁切、倒圆及止水沟加工应先于工厂内加工完毕，不得于现场加工。
b) 安装时各邻接的台面应在同一水平面上，边缘亦应切齐，必须提供适当的支撑装置以避免台面板因支撑力不足而造成塌埳。接合台面板时，接合处的缝隙不得超过3.175mm (1/8”)，并应保持平整。
c) 安装工作完成后，必须小心清除台面上之杂物，并以专用保养油做清洁维护工作。
d) 对于必要的水槽及公用系统配件提供正确位置及大小的开孔。
2.3 水槽:
a) 采台下盆式安装法
b) 安装时台面配合水槽尺寸长宽各缩小30mm开孔，台面开口上缘需采4〜6mm 圆弧R角倒圆处理。
c) 水槽安装时应以钢制悬吊托架(采悬挂式托架，不得将托架直接置于柜体底板处向上支撑，以免占据其下之置物空间)由底部向上承托支撑至与台面底部接合，其钢制托架强度需能支持水槽满水位之重量而不下垂变形，水槽上缘与台面底部齐平，其与台面接缝处采用中性防霉硅胶填充，以达到一体及防漏的效果。
d) 所有水槽应配置有防虹吸瓶式回收器。
2.4 所有实验桌底柜均应为独立完整之单元，左右双开门柜体需为开放式设计，不得设置中间支撑柱以免影响物品存取，落地型底柜需附钢制调整脚，底柜下缘距地面高度应大于10mm。
2.5 所有钢制柜体及框架均须经良好且完整的耐酸碱EPOXY粉末烤漆工序处理，成品烤漆膜厚平均值75μm及以上。
2.6 实验桌各单元底柜之门片及抽屉应易于开关，而不会有任何阻碍。
2.7 检查各相邻的底柜同一型式及高度的门及抽屉应在同一中线而不可因加工不良而造成高低参差不齐，其中线的差距应在3mm(1/8”)内。
2.8 检查实验桌外部所有的开孔，其上均应用适当之孔塞塞住。
2.9 确认所有实验桌环氧树脂台面是否以环氧树脂专用接合剂而非硅胶接合，为求美观档水背板与台面的接合可使用硅胶。各接合处的缝隙宽度不可超过3mm (1/8”)，接合后的台面应平整而不可高低不齐。
2.10 实验桌台面板应尽量保持其完整性，亦应尽量减少接合，且不可以纵向方式接和。
2.11 检查实验桌台面板靠外侧的边缘下方是否有止水沟，止水沟的加工方式为内 凹式而非隆起式，止水沟宽度应不少于2mm，深度不少于3mm。
2.12 检查水、电、气及风管的位置和接口，检查是否与设计有偏差，及时将信息 反馈到项目经理，由项目经理与业主或其它承包商进行沟通，确定最后解决方案。
2.13 将底柜按照图纸进行排列，同一实验台的所有柜子在安装位置将底柜的上沿调在同一水平面上（绝对禁止在没有调整水平的情况下加垫木片等物品）。水平调整完毕后，要求底柜与底柜的缝隙距离不超过1mm，接缝的上下误差不超过0.5mm。如果需要将所有的底柜连成一体，必须在调完水平后进行连接，连接时必须要求接缝处均匀紧密接缝之间的距离不超过0.7mm，上下接缝误差大小不能超过0.5mm。
2.14 柜子调完水平后，进行管道安装
2.15 管道安装完毕后，将台面置于底柜上，台面与台面的接缝距离不能超过 2mm，接缝两头距离误差大小不能超过0.5mm，检查台面是否弯曲变形，如果台面有细微的变形，在空隙处加入橡胶垫片，确保台面安装稳固。台面安装时，如果台面有保护膜，必须在安装完毕后,验收前,才除去保护膜.
2.16 工作时如果人必须踩台面，必须做好台面及家具保护工作.
2.17 冷热水化验龙头：采用美国进口铜制立式水龙头，表面经耐酸碱EPOXY粉 末烤漆处理。出水口2只，上方一口(冷热水混合出口)可360°旋转，下方一口 (冷水)为固定式，两口左右距离 220mm(含)以上，并均具可拆式高压缩口出水口以方便衔接软管
2.18 精密仪器用气体球阀：采用欧美进口不锈钢制双口气体球阀接口，出口管径。
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QGIS API Documentation 3.2.0-Bonn (bc43194)
qgsmaplayerdependency.h
Go to the documentation of this file.
1
2 /***************************************************************************
3 qgsmaplayerdependency.h - description
4 -------------------
5 begin : September 2016
6 copyright : (C) 2016 by Hugo Mercier
7 email : hugo dot mercier at oslandia dot com
8 ***************************************************************************/
9
10 /***************************************************************************
11 * *
12 * This program is free software; you can redistribute it and/or modify *
13 * it under the terms of the GNU General Public License as published by *
14 * the Free Software Foundation; either version 2 of the License, or *
15 * (at your option) any later version. *
16 * *
17 ***************************************************************************/
18
19 #ifndef QGSMAPLAYERDEPENDENCY_H
20 #define QGSMAPLAYERDEPENDENCY_H
21
22 #include "qgis_core.h"
23 #include <QString>
24
35 class CORE_EXPORT QgsMapLayerDependency
36 {
37 public:
39 enum Type
40 {
41 PresenceDependency = 1, //< The layer must be already present (in the registry) for this dependency to be resolved
42 DataDependency = 2 //< The layer may be invalidated by data changes on another layer
43 };
44
46 enum Origin
47 {
48 FromProvider = 0, //< Dependency given by the provider, the user cannot change it
49 FromUser = 1 //< Dependency given by the user
50 };
51
53 QgsMapLayerDependency( const QString &layerId, Type type = DataDependency, Origin origin = FromUser )
54 : mType( type )
55 , mOrigin( origin )
56 , mLayerId( layerId )
57 {}
58
60 Type type() const { return mType; }
61
63 Origin origin() const { return mOrigin; }
64
66 QString layerId() const { return mLayerId; }
67
69 bool operator==( const QgsMapLayerDependency &other ) const
70 {
71 return layerId() == other.layerId() && origin() == other.origin() && type() == other.type();
72 }
73
74 #ifdef SIP_RUN
75 long __hash__() const;
77 % MethodCode
78 sipRes = qHash( *sipCpp );
79 % End
80 #endif
81 private:
82 Type mType;
83 Origin mOrigin;
84 QString mLayerId;
85 };
86
87 #ifndef SIP_RUN
88
92 inline uint qHash( const QgsMapLayerDependency &dep )
93 {
94 return qHash( dep.layerId() ) + dep.origin() + dep.type();
95 }
96 #endif
97
98 #endif
Origin origin() const
Returns the dependency origin.
bool operator==(const QgsMapLayerDependency &other) const
Comparison operator.
Type type() const
Returns the dependency type.
Type
Type of dependency.
uint qHash(const QgsMapLayerDependency &dep)
global qHash function for QgsMapLayerDependency, so that it can be used in a QSet ...
QString layerId() const
Returns the ID of the layer this dependency depends on.
QgsMapLayerDependency(const QString &layerId, Type type=DataDependency, Origin origin=FromUser)
Standard constructor.
This class models dependencies with or between map layers.
Origin
Origin of the dependency.
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Posture Correction Shirt
4,407
78
12
Published
Introduction: Posture Correction Shirt
MetaSkin is a technology that is based on the amount of Infrared light which is being reflected by stretchable cloth. With the technology used in MetaSkin we want to develop a way to improve the posture of people during sitting for example. When a person sits straight with the shoulders turned backwards and the breast to the front the cloth on the back of the person just underneath the shoulder blades won’t be stretched. If the person sits crooked or slumped the cloth on the back of the person will be stretched. This difference should evoke a buzzer, sound or vibration to give feedback to the user about his or her posture.
Link to video-https://www.youtube.com/watch?v=-2Lryl2kwdc
Step 1: How Does It Work?
A Reflective sensor which contains a emitter and a receiver part is placed against the cloth. When the cloth is not stretched, a lot of IR light which is sent by the emitter is received again by the receiver. When the cloth is stretched there is less material in front of the sensor which means that less IR light will be reflected by the cloth and received by the receiver. In our product this difference will be used to give warnings to the user. This difference should evoke a buzzer, sound or vibration to give feedback to the user about his or her posture.
Step 2: Do You Crook or Slump While Working on Computer a Lot?
Did you know that experts estimate that as many as 80% of the population will experience a back problem at some time in our lives (American Chiropractic Association, 2015). Most of the people are aware of this but are not able to take action while working.
Step 3: What If You Have a Persuasive Trigger to Improve Your Posture?
Step 4: Materials Required
1. Jumper wires
2. Clothing
3. Long wires
4. LED
5. IR sensor (emitter and receiver)
6. Resistor (10k)
7. Breadboard
8. Arduino
9. Glue
10. Sensor mount
Step 5: Make the Connections
Step 6: Glue the Sensor on Clothing
Make sure the sensor is not glued and fabric is able to stretch under the sensor. Glue the sensor mount or the wires to the fabric
Step 7: Upload the Code to the Arduino Board
Please find attached the code in the pdf attached.
Step 8: Check and Adjust the 'posture' Threshold
Move your posture from straight to bending position and check the change of values of 'posture' in seek monitor (accessible from magnifying icon on top right corner of Arduino interface). Change the value in the code to the lowest average value you see when you are sitting in bending position.
Step 9: Future Possibilities
You can make it wireless and use flexible sensors to make it less clumsy.
Project by-
Renate Hulst and Vinay Bhajantri
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Please be positive and constructive.
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12 Comments
Hey.. I would like to knw tht what kind of arduino is used in this project.. Plus what kind of jumper wires nd sensor mount is used?.
Yes, you can use a flex sensor may be. We just went for a cheaper alternative i.e. is the IR sensor (it costs you USD 1 and flex sensor is around USD 8).
But agreed, it may be a better option!
i may actually make this. i have a flex sensor and the mentioned wifi chip etc. i will let you know.
:) just saw your last step - great concept.!!
Seems like a mild electrical shock would be more effective than a buzzer. ?
Hey! we used- Vishay TCRT1000, Through Hole Reflective Optical Sensor. The link for that sensor is here-http://nl.rs-online.com/web/p/products/7082737/
Awesome concept, I'd like one of these shirts!
Love this!
Just reading this Instructable was enough make me change how I was sitting.
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awoo's blog
By awoo, history, 3 years ago, translation, In English
1279A - New Year Garland
Idea: BledDest
Tutorial
Solution (Roms)
1279B - Verse For Santa
Idea: Roms
Tutorial
Solution (Roms)
1279C - Stack of Presents
Idea: Roms
Tutorial
Solution (Roms)
1279D - Santa's Bot
Idea: BledDest
Tutorial
Solution (Ne0n25)
1279E - New Year Permutations
Idea: Neon
Tutorial
Solution (Ne0n25)
1279F - New Year and Handle Change
Idea: vovuh
Tutorial
Solution (vovuh)
• Vote: I like it
• +88
• Vote: I do not like it
| Write comment?
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3 years ago, # |
Vote: I like it 0 Vote: I do not like it
Can someone guide me through PROBLEM D.Im not getting the same result as the test case.Can someone simplify the tutorial.
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3 years ago, # ^ |
Rev. 2 Vote: I like it 0 Vote: I do not like it
Probability of choosing 1 kid in among n is 1/n. Let that kid is "x".
Probability of choosing 1 item from the list of the chosen kid(x) is 1/k. Let that item is "y".
So probability of choosing a pair (x, y) is 1/n AND 1/k.
Now lets choose 1 kid from n kids again. And probability of choosing this kid is 1/n (it is not 1/(n-1) because it is independent form the previous choice. Let this kid is "z".
As there are n kids and all the items of a kids list are different so we can say the probability of choosing a kid who want y item is frequency(y)/n.
So probability of choosing a z kid who want y item is 1/n AND frequency(y)/n.
Probability of finding a valid triple (x, y, z) is (1/n AND 1/k) AND (1/n AND frequency(y)/n).
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3 years ago, # ^ |
Vote: I like it 0 Vote: I do not like it
then what...?? do we have to sum all the probabilities of different students..?? sorry but i didn't get that. pls explain by the first test case. and i don't know why i am getting wrong answer for test case 1 but for case 2 answer is perfect.
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3 years ago, # |
Vote: I like it 0 Vote: I do not like it
Query Contest
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3 years ago, # |
Vote: I like it +8 Vote: I do not like it
In the solution to F, why is it necessary to include the $$$res$$$ variable in the binary search? It seems odd since the last value of $$$res$$$ would always be a $$$mid$$$ where $$$check(mid)$$$ returned more than the permitted operations. In which case the later $$$check(res)<=k$$$ would not trigger. However, replacing it with $$$hi + 1$$$ gives WA on 78.
Otherwise, thanks for a thorough explanation!
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3 years ago, # |
Vote: I like it 0 Vote: I do not like it
Can someone explain the "standard method of lexicographic recovery" of problem E more in detail? I never heard of a standard method for that before.
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3 years ago, # ^ |
Vote: I like it +1 Vote: I do not like it
Similar to the method to solve "restore the k-th permutation". I think you can google the solution to that. Just instead of choosing the next number to put we choose the entire block of numbers.
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3 years ago, # ^ |
Vote: I like it +1 Vote: I do not like it
Try values in the prefix and have a count function to calculate the number of different suffixed having some property. Something like:
code
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3 years ago, # ^ |
Rev. 3 Vote: I like it +11 Vote: I do not like it
Thanks. Now I think I have a better idea of the algorithm. It would be something like this:
First we need to figure out the size of the first block. We try sizes 1, 2, 3, ..., N in that order (because of lexicographic order). We keep a running count of how many permutations we have skipped (initially 0). The number of permutations where the first block has size x is $$$(x-2)! \times DP(x+1)$$$. So we keep skipping block sizes until we find the block size x where the k-th permutation lives. So let's call $$$k' = k - skipped$$$. This means we want the k'-th permutation among all permutations where the first block has size x.
So now we know the first block has size x, but we still don't know the exact permutation of numbers from 1 to x of the first block. To find it, we need to remember that there are (x-2)! good permutations for the first block, and that for each one we can fill the remaining numbers to the right in DP(x+1) ways. This means we need to find the r-th permutation of the first block, where r = k'/DP(x+1). So how do we find the r-th good permutation of a block of size x? Here the idea would be doing something like lexicographic recovery: first in a block of size x, we know the first number is x. So we need to fill the second, third, ..., and x-th numbers of the block. For the second number we can try 1, 3, 4, ..., x-1 (x is already used and 2 would create a self-loop). So this is the part I haven't figured out yet (and I would appreciate some help):
1) How do we check that the prefix of numbers we have placed so far in the block is valid? I guess the naive way of following pointers and making sure there are no closed cycles should be enough (if we close a cycle prematurely then it's wrong).
2) This part I don't have any clue yet: How do we count the number of ways we can fill the remaining numbers to the right (the suffix of the block we haven't filled yet) such that the full permutation is good? We have to count making sure we don't accidentally count permutations that would create invalid cycles in combination with the prefix we have placed so far. No clue how to do that.
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3 years ago, # ^ |
Rev. 2 Vote: I like it 0 Vote: I do not like it
if you have prefix of length $$$X$$$ and set first $$$Y$$$ numbers correctly(no cycle there), the rest can be set in $$$(X - Y - 1)!$$$ ways. You can see it in code as well, where $$$cycl[lft]$$$ is used. You can think about it like this, you have to put $$$X-Y$$$ numbers in rest of the places. It's obvious that no number can go to it's place like $$$p[a] = a$$$ as it will create a cycle of length one. Now putting some number in some place, let's say $$$p[2]=6$$$ creates a constraint $$$p[6]\neq2$$$, but there already was a constraint $$$p[6]\neq6$$$ so for one place there is still one constraint about numbers that are left. So if you know that block has size $$$6$$$ it should start with $$$6$$$ and number of ways is $$$4!$$$. When you move to next position, you basically have same problem, you "know" (by iterating) the first number and size is one less.
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3 years ago, # |
Rev. 2 Vote: I like it +134 Vote: I do not like it
"P. S.: We don't have the strict proof that the ans(c) is convex, but we have faith and stress. We'd appreciate it if someone would share the proof in the comment section."
After discussing the problem with my brother we have come up with a proof of ans(c) being convex. The way I think of the problem is that you are given a list of numbers containing $$$0$$$s and $$$1$$$s, and then let $$$g(k)$$$ be the maximum number of $$$1$$$s you can cover with $$$k$$$ intervals of length $$$L$$$. In this formulation the goal will be to prove that $$$g(k)$$$ is concave.
Take an optimal $$$k$$$ interval solution and an optimal $$$k + 2$$$ interval solution. Then make a bipartite graph out of the intervals, where edges are drawn between intervals that intersect. Here is an example for $$$k=7$$$ (the black bars symbolizes the intervals of the two solutions)
The bipartite graph then becomes
Bipartite graph
The reason for picking $$$k$$$ and $$$k + 2$$$ is because I want to create a $$$k + 1$$$ configuration from this that is at least as good as the mean of $$$g(k)$$$ and $$$g(k + 2)$$$. The method I apply is based on analysis of alternating paths in the bipartite graph, in particular what happens if you "flip" a path (meaning you exchange intervals between the two solutions). For example flipping the 3rd path from the left will result in
The 3rd path from the left has been flipped
The following is a proof of the existence of a $$$k + 1$$$ interval solution that contains at least as many $$$1$$$s as the mean of the optimal $$$k$$$ solution and the optimal $$$k + 2$$$ solution.
Formal proof
So conclusion from this is that $$$g(k + 1) \ge \frac{g(k) + g(k + 2)}{2}$$$, i.e. $$$g$$$ is concave.
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3 years ago, # ^ |
Vote: I like it +29 Vote: I do not like it
Awesome! This kind of proof resembles the proof of convexity for Monge dynamic programming. This also has an additional nice property of being constructive.
For the proofs using less brain, it seems the problem can be reduced into negative cycle canceling. Construct an edge $$$i \rightarrow i + 1$$$ with cost 0, and $$$j \le i$$$ with cost $$$-sum[i, j)$$$ for $$$i < j \le i + L$$$. Since negative cycle canceling is dual to the augmenting path, it should have the same convex property.
In the contest, I thought the problem had an MCMF modeling. But after the contest, what I had in mind turned out to be wrong, and this is the best I could get afterwards.
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3 years ago, # |
Vote: I like it +3 Vote: I do not like it
Just curious as to why is it that some problem names in the tutorial appear in Russian (like problem F), while others are in English? I have seen this in many previous editorials as well.
»
3 years ago, # |
Vote: I like it 0 Vote: I do not like it
For the solution of Verse for Santa, alot of macros are used. Its not that clear to me as i am not that fluent in c++ and macros are not defined in solution. Can someone please give reference to their full solution link. It will help me thanks. If its java then better.
• »
»
3 years ago, # ^ |
Vote: I like it +10 Vote: I do not like it
Its a python solution, not a C++ solution.
»
3 years ago, # |
Vote: I like it 0 Vote: I do not like it
In the first question "New Year garland" For the second test case where 1 red, 2 green and 10 blue bulbs are given. we can arrange them as GRBGBBBBBBBBB But according to the question it should print NO
• »
»
3 years ago, # ^ |
Vote: I like it 0 Vote: I do not like it
No, you can't. It is given that two same lights can't be adjacent.
»
3 years ago, # |
Rev. 2 Vote: I like it 0 Vote: I do not like it
In problem C its not clear to me for case.
7 2
2 1 7 3 4 5 6
3 1
how is output 8. I am having doubts can someone please clarify. Thanks
• »
»
3 years ago, # ^ |
Vote: I like it 0 Vote: I do not like it
First in array A {2 1 7 3 4 5 6},we need to remove first 4 element to grab 3 and putting back the removed elements in order of ({1,2,7} or {1,7,2}) and your array A become {1,2,7,4,5,6}; the total opearations you did before was 7{removing 4 elements + putting back 3 elements} and now you just need 1 more operation to remove 1 from A; so total no of operations = 7+1 = 8;
»
3 years ago, # |
Vote: I like it 0 Vote: I do not like it
I can't understand what (this choice is independent of choosing x and y) means in D problem, my solution works as follow:
1. Find the number of possible triples(x,y,z).
2. Find the number of valid triples(x,y,z).
3. Calculate the probability like that => prob = valid/possible
But that don't give me the correct answer. Can anyone help me?
»
3 years ago, # |
Vote: I like it +8 Vote: I do not like it
I was trying to figure out all possible triplets in Problem D . But I couldn't figure out all 8 possibilities .
Here is what I got ..
for the first case which is ---
2
2 2 1
1 1
(1 , 2 , 1) , (1 , 2 , 2)
(1 , 1 , 1) , (1 , 1 , 2)
(2 , 1 , 1) , (2 , 1 , 2)
Where (x , y , z) is ...
x = kid number .
y = one of the gifts kid x wants
z = the chosen gift y is given to kid z .
Now I am wondering what are 2 other possibilities ?
Can anybody help me ??
• »
»
3 years ago, # ^ |
Vote: I like it +3 Vote: I do not like it
You are right, there are 6 possibilities. But they do not have the same probability.
After first selection, both x = 1 and x = 2 have probability of 1/2. If he selected x = 2, then we have only 1 selection for y, so selection (x = 2, y = 1) also has probability of 1/2. Finally, each of (2, 1, 1) and (2, 1, 2) would be 1/4.
In case of x = 1, both y and z have two options, so probability for each of (1, 2, 1), (1, 2, 2), (1, 1, 1), (1, 1, 2) is 1/8.
Then, since the only incorrect choice is (1, 2, 2), whose probability is 1/8, the answer will be 7/8.
• »
»
»
3 years ago, # ^ |
Vote: I like it 0 Vote: I do not like it
Ouch !
I feel so dumb now . How can I miss that for each kid the sum of probability distribution will be $$$1/n$$$ .
»
3 years ago, # |
Rev. 2 Vote: I like it +16 Vote: I do not like it
In problem F, the claim is $$$ dp[n][k] = res(\lambda_{opt},c_{\lambda_{opt}}) - \lambda_{opt}k $$$.
Why is the last value $$$k$$$ instead of $$$ c_{\lambda_{opt}} $$$?
• »
»
3 years ago, # ^ |
Vote: I like it 0 Vote: I do not like it
You missed the previous paragraph where we claim that $$$res(\lambda_{opt},k)=res(\lambda_{opt},c_{\lambda_{opt}})$$$.
|
__label__pos
| 0.830471 |
v0.36-branch
Search
K
Adding a custom batch materialization engine
Overview
Feast batch materialization operations (materialize and materialize-incremental) execute through a BatchMaterializationEngine.
Custom batch materialization engines allow Feast users to extend Feast to customize the materialization process. Examples include:
• Setting up custom materialization-specific infrastructure during feast apply (e.g. setting up Spark clusters or Lambda Functions)
• Launching custom batch ingestion (materialization) jobs (Spark, Beam, AWS Lambda)
• Tearing down custom materialization-specific infrastructure during feast teardown (e.g. tearing down Spark clusters, or deleting Lambda Functions)
Feast comes with built-in materialization engines, e.g, LocalMaterializationEngine, and an experimental LambdaMaterializationEngine. However, users can develop their own materialization engines by creating a class that implements the contract in the BatchMaterializationEngine class.
Guide
The fastest way to add custom logic to Feast is to extend an existing materialization engine. The most generic engine is the LocalMaterializationEngine which contains no cloud-specific logic. The guide that follows will extend the LocalProvider with operations that print text to the console. It is up to you as a developer to add your custom code to the engine methods, but the guide below will provide the necessary scaffolding to get you started.
Step 1: Define an Engine class
The first step is to define a custom materialization engine class. We've created the MyCustomEngine below.
from typing import Any, Callable, Dict, List, Optional, Sequence, Tuple, Union
from feast.entity import Entity
from feast.feature_view import FeatureView
from feast.batch_feature_view import BatchFeatureView
from feast.stream_feature_view import StreamFeatureView
from feast.infra.materialization import LocalMaterializationEngine, LocalMaterializationJob, MaterializationTask
from feast.infra.offline_stores.offline_store import OfflineStore
from feast.infra.online_stores.online_store import OnlineStore
from feast.repo_config import RepoConfig
class MyCustomEngine(LocalMaterializationEngine):
def __init__(
self,
*,
repo_config: RepoConfig,
offline_store: OfflineStore,
online_store: OnlineStore,
**kwargs,
):
super().__init__(
repo_config=repo_config,
offline_store=offline_store,
online_store=online_store,
**kwargs,
)
def update(
self,
project: str,
views_to_delete: Sequence[
Union[BatchFeatureView, StreamFeatureView, FeatureView]
],
views_to_keep: Sequence[
Union[BatchFeatureView, StreamFeatureView, FeatureView]
],
entities_to_delete: Sequence[Entity],
entities_to_keep: Sequence[Entity],
):
print("Creating new infrastructure is easy here!")
pass
def materialize(
self, registry, tasks: List[MaterializationTask]
) -> List[LocalMaterializationJob]:
print("Launching custom batch jobs or multithreading things is pretty easy...")
return [
self._materialize_one(
registry,
task.feature_view,
task.start_time,
task.end_time,
task.project,
task.tqdm_builder,
)
for task in tasks
]
Notice how in the above engine we have only overwritten two of the methods on the LocalMaterializatinEngine, namely update and materialize. These two methods are convenient to replace if you are planning to launch custom batch jobs.
Step 2: Configuring Feast to use the engine
Configure your feature_store.yaml file to point to your new engine class:
project: repo
registry: registry.db
batch_engine: feast_custom_engine.MyCustomEngine
online_store:
type: sqlite
path: online_store.db
offline_store:
type: file
Notice how the batch_engine field above points to the module and class where your engine can be found.
Step 3: Using the engine
Now you should be able to use your engine by running a Feast command:
feast apply
Registered entity driver_id
Registered feature view driver_hourly_stats
Deploying infrastructure for driver_hourly_stats
Creating new infrastructure is easy here!
It may also be necessary to add the module root path to your PYTHONPATH as follows:
PYTHONPATH=$PYTHONPATH:/home/my_user/my_custom_engine feast apply
That's it. You should now have a fully functional custom engine!
|
__label__pos
| 0.848537 |
CYP-Specific PBPK/PD Models to Interpret Biomarkers for Organophosphate Pesticides
EPA Grant Number: R833454
Title: CYP-Specific PBPK/PD Models to Interpret Biomarkers for Organophosphate Pesticides
Investigators: Olson, James , Bonner, Matthew R. , Browne, Richard , Knaak, James B. , Kostyniak, Paul J. , Yu, Aiming
Institution: The State University of New York at Buffalo
EPA Project Officer: Hahn, Intaek
Project Period: September 1, 2007 through August 31, 2010
Project Amount: $749,612
RFA: Interpretation of Biomarkers Using Physiologically Based Pharmacokinetic Modeling (2006) RFA Text | Recipients Lists
Research Category: Health Effects , Health
Objective:
The primary objective of the proposed studies is to improve existing models to better estimate exposures, target tissue dose and resulting effects in human populations, utilizing the abundance of urinary metabolite / biomarker data for the organophosphate (OP) pesticides, chlorpyrifos, parathion, methyl parathion, and diazinon. Current physiologically based pharmacokinetic /pharmacodynamic (PBPK/PD) models, which utilize single kinetic constants (Km, Vmax) for a given metabolic pathway, using rat liver microsomes, do not accurately reflect human enzymes or the variability in content and activity of key enzymes (cytochrome P450s, CYPs; paraoxonase-1, PON1) which regulate the metabolic activation and detoxification of OPs in human populations exposed to pesticides. It is hypothesized that more accurate measures of exposure, target tissue dose and subsequent effects will come from existing PBPK/PD models, which incorporate human CYP-specific kinetic parameters (Km and Vmax) for OP metabolism, CYP–specific content in the liver, and the function and content of genetic variants in key enzymes (CYP2B6, CYP2C19, PON1) which regulate OP metabolic activation and detoxification.
Approach:
1. Generate CYP-specific PBPK/PD models, utilizing existing urinary metabolite / biomarker data for chlorpyrifos, parathion, methyl parathion, and diazinon, to better estimate exposure, target tissue dose and effects of these OPs in human populations by incorporating 1) kinetic parameters (Km and Vmax) for OP metabolism by specific human CYPs, 2) hepatic content of specific CYPs, 3) PON1 activity and content, and 4) serum protein binding data for OPs and metabolites into existing PBPK/PD models.
2. Generate the kinetic parameters for OP metabolism by prevalent polymorphic isoforms of CYP2B6, CYP2C19 and PON1, which are known to exhibit altered enzyme activity and/or expression (biomarkers of susceptibility).
3. Integrate the altered activity and expression data from the key polymorphic CYPs and PON1 (aim 2) into the CYP-specific PBPK/PD models (aim 1).
4. Incorporate the data from aims 1-3 into a multiroute, multichemical, PBPK/PD model (ERDEM, Exposure Related Dose Estimating Model) that can convert urine biomarker data for several OP pesticides to measures of exposure, target tissue dose and effect (inhibition of acetylcholinesterase) in populations of interest.
5. Expected Results:
Since the hepatic levels of specific CYPs exhibit marked variability across both population and age groups, the proposed CYP-specific PBPK/PD models should prove to be more accurate and more easily modified to address factors such as age, CYP content and genetic polymorphisms in CYPs and PON1 (biomarkers of susceptibility). CYP-Specific PBPK/PK models for OP pesticides, that better estimate population variability in exposure, target tissue dose and effect, will be valuable tools for risk assessment efforts for these OPs in the general population, special exposure groups, susceptible individuals (polymorphisms in CYPs and /or PON1), and 6-11 year old children, which are the age group in the U.S. with the highest levels of urinary metabolites specific for these commonly used OPs (CDC 2005).
Supplemental Keywords:
human health, sensitive populations, infants, esterase, oxon, trichloro-2-pyridinol (TCPy), p-nitrophenol (PNP), 2-isopropyl-4-methyl-6-hydroxypyrimidine,, RFA, Health, Scientific Discipline, Health Risk Assessment, Endocrine Disruptors - Environmental Exposure & Risk, endocrine disruptors, Risk Assessments, Biochemistry, Biology, Endocrine Disruptors - Human Health, bioindicator, neurotoxic, dose response, exposure studies, thyroid toxicants, developmental biology, human growth and development, toxicity, hormone production, ecological risk assessment model
Progress and Final Reports:
6. 2008
7. 2009
8. Final
|
__label__pos
| 0.594863 |
Gabapentin Withdrawal & Detox
Gabapentin is a generic drug that is also available under the brand name Neurontin. Gabapentin is used primarily for the treatment of partial seizures in people with epilepsy. It’s also used to treat nerve pain related to shingles infections and diabetes. There are off-label uses of gabapentin as well. For example, it’s sometimes used off-label to treat restless leg syndrome and anxiety disorders. Gabapentin is an analog of the natural neurotransmitter GABA. When taken, gabapentin acts as a central nervous system depressant. Some of the common side effects of the drug include drowsiness and dizziness. Severe side effects may occur as well -but are rare. For example, some people may experience suicidal thoughts or behaviors, or changes in their mood or behavior.
While gabapentin does have therapeutic benefits, especially for people who have epilepsy, there is a risk of abuse. Gabapentin abuse has become a larger problem in recent years. Gabapentin isn’t a controlled substance, making it easier to get than drugs like narcotics and benzodiazepines. When the drug is used and abused, both addition and dependence are possible. If someone is dependent upon gabapentin and they stop using it suddenly or lower their dose, they may experience withdrawal symptoms.
What Causes Gabapentin Withdrawal?
Gabapentin crosses the blood-brain barrier and interacts with GABA receptors. Any time a drug affects the brain, the brain starts to change how it functions in response. For example, if the brain becomes used to the presence of gabapentin, it may produce less of its own natural GABA. Over time, the brain adapts to the presence of these drugs. When someone suddenly stops using a drug, the brain goes into a type of shock and struggles to regain a sense of normalcy in its function. That struggle is why withdrawal symptoms occur. Any time a drug affects GABA, withdrawal symptoms can be particularly severe. For example, because both alcohol and benzodiazepines affect GABA, they have some of the most dangerous withdrawal symptoms.
Withdrawal Factors
Not everyone who uses gabapentin is going to go through withdrawal. Even if someone does go through withdrawal from gabapentin, it doesn’t always have to be severe. Factors that influence the severity of withdrawal symptoms include the amount of time someone has used the drug, the dosage they usually take, and whether or not other substances are also involved. If someone has been taking a high dose of gabapentin every day for several years, they’re probably going to have more severe withdrawal symptoms than someone who has been taking it as prescribed for a few months. If a doctor works with a patient to gradually taper down their dosage of gabapentin over time, as opposed to stopping suddenly, that individual may not go through a noticeable withdrawal at all.
Gabapentin Withdrawal & Detox
Some of the most common symptoms of gabapentin withdrawal include sweating, dizziness and fatigue. People may experience headaches, insomnia and other sleep disturbances, and muscle pain as well. It’s common for people to have anxiety when they go through withdrawal from gabapentin, as well as depression. Spasms, restlessness and stomach pain can occur. Some of the most serious possible symptoms of gabapentin withdrawal include seizures and suicidal thoughts or behaviors.
For most people, gabapentin withdrawal will start anywhere from 12 to 48 hours after they’ve taken the last dose of the drug. In the initial hours of gabapentin withdrawal, symptoms can be the most serious. These early symptoms, from the first 12 to 24 hours, can include anxiety, sweating, headaches, tremors and seizures. By day three, most people will have gone through the worst of the gabapentin withdrawal symptoms, such as restlessness, confusion and agitation. Some people may have more serious symptoms like hallucinations and changes in the heart rate. Within around a week, the majority of the gabapentin withdrawal symptoms should have subsided. Some people may have post-acute withdrawal symptoms, or PAWs, however. These are ongoing symptoms that are primarily psychological, such as anxiety or depression.
Gabapentin Withdrawal & Detox
Some people may be able to go through detox for gabapentin withdrawal at home by following their doctor’s instructions. For the most part, however, detox for gabapentin withdrawal is best done in a professional environment. With any drug that affects GABA, because there is a potential for such severe side effects, a medical detox is usually the safest option. During a medical detox, a patient can be closely monitored and provided with the necessary treatments to keep them safe and comfortable.
People often wonder what to expect from a detox center -whether it’s for gabapentin or another substance. When someone goes to a detox center, they will typically answer assessment questions and move through the intake process. This is when the staff can assess the individual situation and make a plan to provide the best care for the patient. This is one of the most important elements of addiction treatment because it can determine and guide the entire treatment process. When a patient is transparent and honest at the beginning of detox, they can receive the optimal detox treatment. The primary objective during detox is to mitigate or alleviate the unpleasant symptoms of withdrawal in order to increase the patient’s likelihood of success.
The Recovery Village works to create customized, tailored detox and treatment plans that can change the course of your life, or the life of your loved one. Contact us today to learn more.
|
__label__pos
| 0.903599 |
RSS
Monthly Archives: September 2012
Restore MySql Backup file by using JAVA
First read my early blog post “Take the MySql DB backup in JAVA” and you can download the source from there
public boolean restoreDatabase(String dbUserName, String dbPassword, String source) {
String[] executeCmd = new String[]{"mysql", "--user=" + dbUserName, "--password=" + dbPassword, "-e", "source " + source};
Process runtimeProcess;
try {
runtimeProcess = Runtime.getRuntime().exec(executeCmd);
int processComplete = runtimeProcess.waitFor();
if (processComplete == 0) {
log.info("Backup restored successfully with " + source);
return true;
} else {
log.info("Could not restore the backup " + source);
}
} catch (Exception ex) {
log.error(ex, ex.getCause());
}
return false;
}
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4 Comments
Posted by on September 28, 2012 in java, mysql
Tags: , ,
Take the MySql DB backup in JAVA
This is very impotent blog post for the most application developers
Here I’m talking about how to backup the mysql databasae through JAVA program.
Download MySqlBackup.java you have to download mysql.jar and log4j.jar
1. Following method will create the data base backup without create and drop database commands
It’s only creating tables and included data
When your going to restore the database first you have to create new DB
public boolean backupDataWithOutDatabase(String dumpExePath, String host, String port, String user, String password, String database, String backupPath) {
boolean status = false;
try {
Process p = null;
DateFormat dateFormat = new SimpleDateFormat("dd-MM-yyyy");
Date date = new Date();
String filepath = "backup(without_DB)-" + database + "-" + host + "-(" + dateFormat.format(date) + ").sql";
String batchCommand = "";
if (password != "") {
//only backup the data not included create database
batchCommand = dumpExePath + " -h " + host + " --port " + port + " -u " + user + " --password=" + password + " " + database + " -r \"" + backupPath + "" + filepath + "\"";
} else {
batchCommand = dumpExePath + " -h " + host + " --port " + port + " -u " + user + " " + database + " -r \"" + backupPath + "" + filepath + "\"";
}
Runtime runtime = Runtime.getRuntime();
p = runtime.exec(batchCommand);
int processComplete = p.waitFor();
if (processComplete == 0) {
status = true;
log.info("Backup created successfully for without DB " + database + " in " + host + ":" + port);
} else {
status = false;
log.info("Could not create the backup for without DB " + database + " in " + host + ":" + port);
}
} catch (IOException ioe) {
log.error(ioe, ioe.getCause());
} catch (Exception e) {
log.error(e, e.getCause());
}
return status;
}
2. This method will create a backup from Specified database with create and drop database commands
When your going to restore the DB you don’t need to create the DB it will handle by the script
public boolean backupDataWithDatabase(String dumpExePath, String host, String port, String user, String password, String database, String backupPath) {
boolean status = false;
try {
Process p = null;
DateFormat dateFormat = new SimpleDateFormat("dd-MM-yyyy");
Date date = new Date();
String filepath = "backup(with_DB)-" + database + "-" + host + "-(" + dateFormat.format(date) + ").sql";
String batchCommand = "";
if (password != "") {
//Backup with database
batchCommand = dumpExePath + " -h " + host + " --port " + port + " -u " + user + " --password=" + password + " --add-drop-database -B " + database + " -r \"" + backupPath + "" + filepath + "\"";
} else {
batchCommand = dumpExePath + " -h " + host + " --port " + port + " -u " + user + " --add-drop-database -B " + database + " -r \"" + backupPath + "" + filepath + "\"";
}
Runtime runtime = Runtime.getRuntime();
p = runtime.exec(batchCommand);
int processComplete = p.waitFor();
if (processComplete == 0) {
status = true;
log.info("Backup created successfully for with DB " + database + " in " + host + ":" + port);
} else {
status = false;
log.info("Could not create the backup for with DB " + database + " in " + host + ":" + port);
}
} catch (IOException ioe) {
log.error(ioe, ioe.getCause());
} catch (Exception e) {
log.error(e, e.getCause());
}
return status;
}
10 Comments
Posted by on September 28, 2012 in java, mysql
Tags: , ,
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Issue Tracking system – https://www.etraxis.com/
2 Comments
Posted by on September 18, 2012 in java, MSSQL, mysql, Other
Tags: , , , ,
How to Create Sample Database in MSSQL with Sample Data
This is very useful for application testing.
First you have to download the relevant sample databases from http://msftdbprodsamples.codeplex.com/releases/view/59211
Then copy the relevant AdventureWorks2008R2_Data.mdf file in to your machine where the MSSQL is installed. then go to MSSQL Editor and run the following query
CREATE DATABASE AdventureWorks2012
ON (FILENAME = 'D:\AdventureWorks2012_Data.mdf')
FOR ATTACH_REBUILD_LOG ;
Don’t delete the above AdventureWorks2008R2_Data.mdf because this will be your database in MSSQL
Leave a comment
Posted by on September 14, 2012 in MSSQL
Tags: , , ,
About log4j.properties
log4j.properties Syntax:
Following is the syntax of log4j.properties file for an appender X:
# Define the root logger with appender X
log4j.rootLogger = DEBUG, X
# Set the appender named X to be a File appender
log4j.appender.X=org.apache.log4j.FileAppender
# Define the layout for X appender
log4j.appender.X.layout=org.apache.log4j.PatternLayout
log4j.appender.X.layout.conversionPattern=%m%n
log4j.rootLogger – the first element should be the logging level, Second one should be the appender
Logging levels –
TRACE < DEBUG < INFO < WARN < ERROR < FATAL < ALL
Level Description
ALL All levels including custom levels.
DEBUG Designates fine-grained informational events that are most useful to debug an application.
ERROR Designates error events that might still allow the application to continue running.
FATAL Designates very severe error events that will presumably lead the application to abort.
INFO Designates informational messages that highlight the progress of the application at coarse-grained level.
OFF The highest possible rank and is intended to turn off logging.
TRACE Designates finer-grained informational events than the DEBUG.
WARN Designates potentially harmful situations.
We have used only one appender FileAppender in our example above. All the possible appender options are:
• AppenderSkeleton
• AsyncAppender
• ConsoleAppender
• DailyRollingFileAppender
• ExternallyRolledFileAppender
• FileAppender
• JDBCAppender
• JMSAppender
• LF5Appender
• NTEventLogAppender
• NullAppender
• RollingFileAppender
• SMTPAppender
• SocketAppender
• SocketHubAppender
• SyslogAppender
• TelnetAppender
• WriterAppender
We have used PatternLayout with our appender. All the possible options are:
• DateLayout
• HTMLLayout
• PatternLayout
• SimpleLayout
• XMLLayout
Leave a comment
Posted by on September 13, 2012 in java
Tags: , , ,
Log4j with NetBeans IDE
I think this post will help you to configure the log4j in Netbeans IDE
First you have to create the Java Project “Log4j”
Then you have to put the “log4j.properties” file in to your src folder “This should be located in root”
#### Use Three appenders,
#stdout - is used for write to console
#R - is used for write to file
log4j.rootLogger=debug, stdout, R
# Print only messages of priority WARN or higher for your category
# log4j.category.your.category.name=WARN
# Specifically inherit the priority level
#log4j.category.your.category.name=INHERITED
# Print only messages of level WARN or above in the package
#This is use for debuging mode
log4j.logger.testlogging=DEBUG
#### Appender writes to console
log4j.appender.stdout=org.apache.log4j.ConsoleAppender
log4j.appender.stdout.layout=org.apache.log4j.PatternLayout
# Pattern to output the caller's file name and line number.
log4j.appender.stdout.layout.ConversionPattern=%d{yyyy-MM-dd} %5p [%t] (%F:%L) - %m%n
#### Appender writes to a file
#log4j.appender.R=org.apache.log4j.FileAppender
log4j.appender.R=org.apache.log4j.RollingFileAppender
log4j.appender.R.File=example.log
# Control the maximum log file size
log4j.appender.R.MaxFileSize=100KB
# Archive log files (one backup file here)
log4j.appender.R.MaxBackupIndex=1
log4j.appender.R.layout=org.apache.log4j.PatternLayout
log4j.appender.R.layout.ConversionPattern=%d{yyyy-MM-dd} %5p [%t] (%F:%L) - %m%n
#log4j.appender.R.layout.ConversionPattern=%n%p - %m
Finally you have to take this “log4j-1.2.17.jar” file and add in to the class path
Now you done all configuration lets try to use it
package log4j;
import org.apache.log4j.Logger;
/**
*
* @author dinuka
*/
public class Log4J {
//initializing the logger
static Logger log = Logger.getLogger(Log4J.class.getName());
/**
* @param args the command line arguments
*/
public static void main(String[] args) {
//logging in different levels
log.trace("This is a Trace");
log.debug("This is a Debug");
log.info("This is an Info");
log.warn("This is a Warn");
log.error("This is an Error");
log.fatal("This is a Fatal");
}
}
Output file
2012-09-11 TRACE [main] (Log4J.java:19) - This is a Trace
2012-09-11 DEBUG [main] (Log4J.java:20) - This is a Debug
2012-09-11 INFO [main] (Log4J.java:21) - This is an Info
2012-09-11 WARN [main] (Log4J.java:22) - This is a Warn
2012-09-11 ERROR [main] (Log4J.java:23) - This is an Error
2012-09-11 FATAL [main] (Log4J.java:24) - This is a Fatal
10 Comments
Posted by on September 11, 2012 in java
Tags: , ,
Read XML in Java
First follow this post https://malalanayake.wordpress.com/2012/09/07/write-xml-in-java/
XmlReader.java
package mysamples;
import java.io.FileInputStream;
import java.io.IOException;
import java.io.InputStream;
import java.util.ArrayList;
import java.util.HashMap;
import java.util.Map;
import javax.xml.parsers.DocumentBuilder;
import javax.xml.parsers.DocumentBuilderFactory;
import javax.xml.parsers.ParserConfigurationException;
import org.w3c.dom.Document;
import org.w3c.dom.Element;
import org.w3c.dom.NamedNodeMap;
import org.w3c.dom.Node;
import org.w3c.dom.NodeList;
import org.xml.sax.SAXException;
public class XmlReader {
private String name;
private String content;
private Map<String, String> nameAttributes = new HashMap<String, String>();
private Map<String, ArrayList<XmlReader>> nameChildren = new HashMap<String, ArrayList<XmlReader>>();
public XmlReader(InputStream inputStream, String rootName) {
this(rootElement(inputStream, rootName));
}
public XmlReader(String filename, String rootName) {
this(fileInputStream(filename), rootName);
}
public XmlReader(String rootName) {
this.name = rootName;
}
private XmlReader(Element element) {
this.name = element.getNodeName();
this.content = element.getTextContent();
NamedNodeMap namedNodeMap = element.getAttributes();
int n = namedNodeMap.getLength();
for (int i = 0; i < n; i++) {
Node node = namedNodeMap.item(i);
String name = node.getNodeName();
addAttribute(name, node.getNodeValue());
}
NodeList nodes = element.getChildNodes();
n = nodes.getLength();
for (int i = 0; i < n; i++) {
Node node = nodes.item(i);
int type = node.getNodeType();
if (type == Node.ELEMENT_NODE) {
XmlReader child = new XmlReader((Element) node);
addChild(node.getNodeName(), child);
}
}
}
public void addAttribute(String name, String value) {
nameAttributes.put(name, value);
}
private void addChild(String name, XmlReader child) {
ArrayList<XmlReader> children = nameChildren.get(name);
if (children == null) {
children = new ArrayList<XmlReader>();
nameChildren.put(name, children);
}
children.add(child);
}
public String name() {
return name;
}
public void setContent(String content) {
this.content = content;
}
public String content() {
return content;
}
public void addChild(XmlReader xml) {
addChild(xml.name(), xml);
}
public void addChildren(XmlReader... xmls) {
for (XmlReader xml : xmls) {
addChild(xml.name(), xml);
}
}
public XmlReader child(String name) {
XmlReader child = optChild(name);
if (child == null) {
throw new RuntimeException("Could not find child node: " + name);
}
return child;
}
public XmlReader optChild(String name) {
ArrayList<XmlReader> children = children(name);
int n = children.size();
if (n > 1) {
throw new RuntimeException("Could not find individual child node: " + name);
}
return n == 0 ? null : children.get(0);
}
public boolean option(String name) {
return optChild(name) != null;
}
public ArrayList<XmlReader> children(String name) {
ArrayList<XmlReader> children = nameChildren.get(name);
return children == null ? new ArrayList<XmlReader>() : children;
}
public String string(String name) {
String value = optString(name);
if (value == null) {
throw new RuntimeException(
"Could not find attribute: " + name + ", in node: " + this.name);
}
return value;
}
public String optString(String name) {
return nameAttributes.get(name);
}
public int integer(String name) {
return Integer.parseInt(string(name));
}
public Integer optInteger(String name) {
String string = optString(name);
return string == null ? null : integer(name);
}
public double doubleValue(String name) {
return Double.parseDouble(optString(name));
}
public Double optDouble(String name) {
String string = optString(name);
return string == null ? null : doubleValue(name);
}
private static Element rootElement(InputStream inputStream, String rootName) {
try {
DocumentBuilderFactory builderFactory = DocumentBuilderFactory.newInstance();
DocumentBuilder builder = builderFactory.newDocumentBuilder();
Document document = builder.parse(inputStream);
Element rootElement = document.getDocumentElement();
if (!rootElement.getNodeName().equals(rootName)) {
throw new RuntimeException("Could not find root node: " + rootName);
}
return rootElement;
} catch (IOException exception) {
throw new RuntimeException(exception);
} catch (ParserConfigurationException exception) {
throw new RuntimeException(exception);
} catch (SAXException exception) {
throw new RuntimeException(exception);
} finally {
if (inputStream != null) {
try {
inputStream.close();
} catch (Exception exception) {
throw new RuntimeException(exception);
}
}
}
}
private static FileInputStream fileInputStream(String filename) {
try {
return new FileInputStream(filename);
} catch (IOException exception) {
throw new RuntimeException(exception);
}
}
}
Here I’m Reading the Employee.xml
public static void main(String[] args) {
XmlReader config = new XmlReader("Employee.xml", "employees");
System.out.println("title: " + config.child("employee").content());
for(XmlReader firstname : config.child("employee").children("firstname")){
System.out.println("First Name :" + firstname.content());
}
for(XmlReader firstname : config.child("employee").children("lastname")){
System.out.println("Last Name :" + firstname.content());
}
for(XmlReader firstname : config.child("employee").children("age")){
System.out.println("Age :" + firstname.content());
}
for(XmlReader firstname : config.child("employee").children("email")){
System.out.println("Email :" + firstname.content());
}
}
If You need to read something like this kind of xml you can do it in this way
<config>
<title>test</title>
<version
major="1"
minor="2"/>
<roles>
<role name="admin"/>
<role name="user"/>
</roles>
<users>
<user name="joe" password="pass" role="admin"/>
<user name="harry" password="secret" role="user" email="[email protected]"/>
</users>
<test/>
</config>
public static void main(String[] args) {
Xml config = new Xml("config.xml","config");
System.out.println("title: "+config.child("title").content());
Xml version = config.child("version");
System.out.println("version: "+version.integer("major")+"."+version.integer("minor"));
for(Xml role:config.child("roles").children("role"))
System.out.println("role: name: "+role.string("name"));
for(Xml user:config.child("users").children("user")) {
String email = user.optString("email");
System.out.println(
"user: name: "+user.string("name")+
", password: "+user.string("password")+
", role: "+user.string("role")+
", email: "+(email==null ? "-" : email));
}
System.out.println("test: "+config.option("test"));
}
Leave a comment
Posted by on September 7, 2012 in java
Tags: , ,
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__label__pos
| 0.965754 |
Flower Whorls and Femaleness
Flower Structure (review)
The flower is basically a strobilus. Here is a generic flower:
The carpel is basically a megasporophyll. It is folded around the ovules to form a hollow chamber, the locule, with walls, the ovary. Since pollen no longer has access to the micropyle of the ovules, it must be deposited on the stigma. The stigma is sticky to hold and promote pollen tube germination. The pollen tube must grow through the style to reach the locule (chamber in the ovary). This growth is fueled by digestion of the female tissues, and by special transfer cells that line the tubular style in many species. Thus, the microgametophyte is dependent upon the female sporophyte tissues for nutrition and logistics. Ultimately the pollen tube finds the micropyle of an ovule, probably by chemotropism. The chemo-attractant may be some combination of Calcium ions, Boron ions, and growth hormones (perhaps auxins).
Meanwhile, in the ovary, the female development occurs. The megasporophyll has the ovule attached to its upper (inner) surface. The site of connection is called the placenta. Xylem and phloem connections cross from the megasporophyll (carpel wall) into the ovule by means of a stalk called the funiculus. Thus the megagametophyte stage will be supported nutritionally by the sporophyte. The table has been turned!
The ovule contains the megasporangium (nucellus) which, in turn, contains one megasporocyte. This diploid cell undergoes meiosis to produce the four haploid megaspores. Typically three of these degenerate, leaving one functional megaspore. The megaspore is not shed and remains in the megasporangium. From here the development varies among the species; what is described next is the "common" pathway--one of several possible developmental sequences found in angiosperms. The megaspore undergoes mitosis to produce 8 haploid nuclei in a free-nuclear megagametophyte. Cytokinesis partitions these eight nuclei into seven cells...the mature megagametophyte.
At the chalazal end (opposite the foot=funiculus) of the megagametophyte are three antipodal cells. These are all that remain of the once-dominant gametophyte thallus! There is the central cell with two nuclei (one from each pole, the so-called polar nuclei). This cell will join with a sperm to make the polyploid endosperm (an autapomorphy of angiosperms). At the micropylar end of the ovule are three cells: the egg cell and two synergids. The egg cell will be joined by a sperm cell to form the zygote. The synergids will enzymatically digest the end of the arriving pollen tube to assist in sperm release. It is possible that the synergids are all that remain of the archegonium sterile jacket cells.
This pathway for ovule development is outlined as follows:
Color Key
sporophytegreen
sporangiumred
sporocyteyellow
sporecyan
gametophyteblue
nucleimagenta
Male nuclei are shown in cyan, the endosperm is shown in yellow.
It is important to note that the details of the embryo sac (the megagametophyte) vary among species of angiosperms. The "Polygonum" type described above applies to the majority of angiosperm species but there are many other patterns. One would study these in a graduate course on angiosperm anatomy.
Syngamy in angiosperms is double, another autapomorphy. One sperm joins with the egg to form the diploid zygote, the other sperm joins with the two polar nuclei in the central cell to form the polyploid endosperm.
The endosperm develops rapidly. Its polyploid genetic status apparently gives it an amazing ability to sequester nutrients from the maternal sporophyte. Sugars are polymerized into starch, amino acids into protein, and acetyl-CoA is maneuvered into lipids. These nutrients will be used later by the developing embryo. Just after karyogamy, the newly formed endosperm cell repeated divides mitotically to form a free-nuclear endosperm cell. This cell can be unbelievably large...the vacuole and liquid cytoplasm of this cell form what we know of as "coconut milk" in the seed of that species of palm for example. If you have had a piña colada, you have drunk liquid endosperm. Later, cytokinesis divides the endosperm into a cellular mass. Indeed the "coconut meat" is solid endosperm that we grate and put on baked goods for example. Solid endosperm of other species (corn, wheat, rye, barley, etc.) is typically ground into flour, mixed into a dough, and then baked to form bread, crackers, pie crust, cake, bagels, etc.
While the endosperm garners nutrients from the maternal sporophyte, the zygote formed by the union of egg and sperm develops. This development involves mitosis and cytokinesis and later morphogenesis. The zygote divides to make a chain of cells called the suspensor. The bulbous cell at the micropylar end of the suspensor is usually connected by a uniseriate filament of cells to a tiny cell at the chalazal end. (Obviously more examples of unequal cytokinesis!) This tiny terminal cell will become the embryo proper. It is pushed into the liquid endosperm region by growth of the suspensor. Is it possible that the suspensor is the evolutionary relic of the haustorial foot in young sporophytes? The suspensor is genetically identical to the embryo proper.
The terminal cell divides mitotically to form a globular embryo...a spherical cluster of a few tiny cells. Polarity is established and the end distal to the suspensor develops two lobes (at least in dicots!) which become leaf primordia and are called cotyledons when mature. The proximal portion of the embryo attached to the suspensor differentiates into a root apex. This stage of development is called the heart stage. The embryo then elongates into what is referred to as the torpedo stage...obviously a vocabulary term from war years in the 20th century. Torpedo means fusiform at one end and with two fin-like projections at the other. This embryo may elongate enough to need to bend within the confines of the ovule, and its development may use up all available endosperm as well. The nutrients in that case are transferred to the embryo...usually stored as reserves in the cotyledons. Thus mature seeds may lack endosperm.
It is important to note that as soon as the included embryo becomes dormant that the name ovule is changed to seed. The seed consists of a seed coat (the ovule integument), storage tissue (endosperm or cotyledon depending on species), and a dormant embryo. This change forces another. The carpel is renamed fruit. As you recall the gymnosperms were noted for their naked seeds. The megasporophyll had a seed on its exposed surface. In angiosperms, because the megasporophyll encloses the seed in a locule, the surrounding ovary wall tissue is the fruit. To disperse seeds then, the carpel must open in some way to shed the seeds within, or it must rely upon an animal to open the carpel to disperse and release the seeds. There will be more on that later. For now, suffice it to say that the fruit is thus an autapomorphic feature of angiosperms. Because this feature is shared by all angiosperms then it is synapomorphic among them.
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| 0.938039 |
/ Images
Image loader component with fallbacks for React Native apps
Image loader component with fallbacks for React Native apps
React Native Image Fallback
Image loader component with fallbacks for React Native apps.
1. What is React Native Image Fallback?
React Native Image Fallback is a lightweight image component which supports fallback images for React Native apps.
2. Getting Started
Install React Native Image Fallback
npm i react-native-image-fallback --save
3. Usage
Import it
import {ImageLoader} from 'react-native-image-fallback';
Use it in you component
render() {
const imageSource = 'http://image.url'; // An image URL
const fallbacks = [
'http://another.image.url', // An image URI
require('./local/image/path'), // A locally require'd image
];
return (
<ImageLoader
source={imageSource}
fallback={fallbacks}
/>
)
}
4. Properties
This is basically a React Native Image. So all the <Image /> props will work. On top of that
• source - The source image. Can be a string URL or a require('') image file
• fallback - The fallback(s). Can be a string URL, a require('') image file or an array consisting of either
• onLoadStart - Accepts a calback function with the first parameter being the image that is being loaded to the component.
• onLoadEnd - Invoked when load either succeeds or fails. Accepts a callback function with the first parameter being the image in question.
• onSuccess - Invoked when the component successfully loads an image. Accepts a callback function with the first parameter being the loaded image.
• onError - Invoked when all the given images fail to load.
GitHub
Comments
|
__label__pos
| 0.986338 |
How to quickly change your branch in Git
As the Editor-in-Chief of TheServerSide, I monitor the search terms that drive readers to the site. One search term that makes me wonder is “git branch change.” This interests me because there are several reasons people could be searching for this topic and come here.
Meaning of ‘git branch change’
If the “git branch change” search term brought you to this site, you are likely interested in one of the following two topics:
1. How do you change and switch between Git branches?
2. How to you change and rename a Git branch?
I’ve written full tutorials on each of those ‘git branch change’ topics, each of which is linked in the numbered list above. But here’s the tl;dr response for each of them.
How to change Git branches
To change a git branch name, you simply switch or checkout the branch of interest and issue the following command:
[email protected] /c/git/github (hotflex)
$ git branch -m hotfix
[email protected] /c/git/github (hotfix)
$ git branch -a
* hotfix
main
Keep in mind that this only performs a local Git branch change. Further steps are required to push the name change to GitHub or GitLab and share the Git branch change with fellow developers.
How to switch between Git branches
Alternatively, a developer who queried “git branch change” may be uninterested in a renaming. Instead, she or he wants to better understand the command to change between Git branches in their local dev environment.
To change Git branches, developers can use either the checkout or switch command. The modern preference is to use git switch. Git switch replaced checkout in a 2020 release, although both commands are still supported.
Git Branch Change
The ability for users to work in isolated branches, and change Git branches when needed, helps simplify distributed version control.
Change Git branches switch command
To change the current Git branch, first list the branches, and then provide the name of the branch of interest to the switch command. In the following example, the user starts on the hotfix branch and then does a git branch change to get on the release branch:
[email protected] /c/git/github (hotfix)
$ git branch -a
* hotfix
main
release
support
development
feature
[email protected] /c/git/github (hotfix)
$ git branch release
[email protected] /c/git/github (release)
$ git status
Whichever path your “git branch change” query takes, be it to switch between branches or rename one, Git makes it relatively easy to perform any of these required functionalities.
App Architecture
Software Quality
Cloud Computing
Security
SearchAWS
Close
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| 0.697236 |
How to use v-if inside v-for loop in vue.js?
by cali_green , in category: Javascript , 6 months ago
How to use v-if inside v-for loop in vue.js?
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1 answer
Member
by deron , 4 months ago
@cali_green
To use v-if inside a v-for loop in Vue.js, you can add the v-if directive on the element within the loop. Here's an example:
1
2
3
4
5
<ul>
<li v-for="item in items" v-if="item.isActive">
{{ item.name }}
</li>
</ul>
In this example, items is an array of objects with name and isActive properties. The v-for directive loops through each item in the array, and the v-if directive conditionally renders the <li> element only if item.isActive is true.
Make sure to place the v-if directive on the same element that you want to conditionally render.
|
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| 0.996435 |
Intereting Posts
Какая ссылка на мою страницу «Архив» по умолчанию? Создайте регистрацию переднего пользователя / login / profile & logout без вмешательства в wp-login.php? Создание программных программ Невозможно загрузить изображения на новую тему Изменение цвета текста виджетов Как сделать копию локальной сборки WordPress в MAMP? Есть ли способ написать собственный HTML-файл в формате .php, а затем вызвать его там, где это необходимо? Правильный способ фильтрации функции get_pages () Получить текущую таксономию и отобразить запрос соответственно Получение сообщения в wp_nav_walker Печать строки в нижний колонтитул с использованием параметра wp_footer Category.php загружает сначала до page.php? Как остановить WooCommerce / WordPress от изменения размера изображения продукта? Как исключить категории из списка get_categories () в панели администратора виджета Обнаружение, если я нахожусь на одной странице без должности, а не на главной странице?
Выполнение произвольного кода по определенному URL-адресу без создания сообщения или страницы?
Я пытаюсь выяснить, какие крючки я могу связать, чтобы прочитать входящий URL-адрес до фактического запуска запроса. Быстрый взгляд говорит мне, может быть, wp_loaded или posts_selection будет хорошим выбором, но будет ненавидеть пропустить что-то, что явно лучше для моих нужд.
template_redirect звучит многообещающе, но, похоже, вызывается после того, как контент выбран, поэтому я не думаю, что это сработает.
Я думал, что код перезаписи может помочь, но подобные generate_rewrite_rules , похоже, не позволяют выполнять произвольный код … кажется, все они направлены на перенаправление на существующий контент.
Моя цель – написать код перенаправления, который мне не нужно привязывать к существующей части содержимого. Я нашел несколько плагинов, которые делают подобные вещи, но все, кажется, требуют, чтобы сообщение или страница были добавлены в систему для перенаправления, к которой нужно подключиться.
Например, я хочу /foo/bar перенаправлять на /baz , но я не хочу создавать страницу foo или подстраницу из foo bar .
@Hi Beau Simensen:
Крючок plugsin_loaded является одним из вариантов, но вы можете использовать плагин, который я опубликовал для этого вопроса, чтобы узнать, какие плагины запущены и в каком порядке и использовать его, чтобы найти подходящий для вас крюк:
• Где я могу найти список крючков WordPress?
Учитывая вашу цель, вы также можете найти этот ответ полезным:
• Создание 301 переадресации для почтовых сообщений, страниц, категорий и URL-адресов изображений?
Если это не то, что вам нужно, проконсультируйтесь, и я надеюсь, что вы предложите больше.
|
__label__pos
| 0.676616 |
Ginseng: Benefits, Uses, and Potential Side Effects
Health Supplements
Ginseng: Benefits, Uses, and Potential Side Effects
Ginseng is a revered herb known for its benefits, diverse uses, and potential side effects in the realm of natural health and Ginseng: Benefits, Uses, and Potential Side Effects. Originating from the roots of plants in the genus Panax, which translates to “all-healing” in Greek, ginseng has been prized for centuries for its purported adaptogenic properties and ability to enhance vitality.
Types of Ginseng
There are several types of ginseng, with the most common being Asian ginseng (Panax ginseng) and American ginseng (Panax quinquefolius). These varieties differ slightly in their chemical composition and therapeutic effects, though both are valued for their medicinal properties.
Active Compounds
The active compounds in ginseng responsible for its health benefits include ginsenosides, polyphenols, and polysaccharides. Ginsenosides, in particular, are believed to contribute to ginseng’s adaptogenic effects, helping the body cope with stress and promoting overall well-being.
Benefits of Ginseng
Ginseng is acclaimed for a range of health benefits, including:
• Boosting energy: It is often used to combat fatigue and enhance physical stamina.
• Supporting cognitive function: Ginseng may improve memory, focus, and mental clarity.
• Enhancing immune function: It helps strengthen the immune system against infections.
• Managing stress: Adaptogenic properties aid in reducing the effects of stress on the body.
• Promoting cardiovascular health: It may lower cholesterol and blood sugar levels.
Uses of Ginseng
Ginseng finds its way into various health supplements, herbal teas, and cosmetic products due to its therapeutic and nourishing qualities. It is also utilized in traditional medicine practices such as Chinese medicine and Ayurveda to treat conditions ranging from low libido to digestive disorders.
Potential Side Effects
While generally considered safe for most people when taken in moderate amounts, ginseng can have side effects including:
• Insomnia and nervousness due to its stimulant effects.
• Digestive issues such as nausea, vomiting, and diarrhea.
• Headaches and dizziness in some individuals.
• Hypoglycemia (low blood sugar) particularly in diabetic patients.
• Allergic reactions like skin rashes or difficulty breathing in rare cases.
Interactions with Medications
Ginseng may interact with certain medications, including blood thinners, diabetes medications, and antidepressants. It’s crucial to consult with a healthcare provider before starting ginseng supplementation, especially if you are on medication or have underlying health conditions.
Formulations and Dosage
Ginseng is available in various forms, including capsules, powders, extracts, and teas. The dosage and formulation can vary based on the intended use and desired benefits. Following recommended guidelines and avoiding excessive consumption is key to minimizing side effects and maximizing efficacy.
Research and Evidence
Scientific research continues to explore the potential of ginseng in various health applications. Studies have shown promising results regarding its anti-fatigue effects, cognitive enhancement, and immune modulation properties. Ongoing research aims to further elucidate its mechanisms of action and expand its therapeutic applications.
Cultivation and Sustainability
Due to high demand, ginseng cultivation has become a significant industry in regions where it thrives, such as Asia and parts of North America. Sustainable practices are crucial to ensure the longevity of ginseng cultivation and preserve its medicinal properties for future generations.
Conclusion
In conclusion, ginseng remains a cornerstone in natural health and traditional medicine, revered for its beneficial effects on energy, cognition, and overall well-being. While it offers numerous health benefits, it’s essential to be aware of potential side effects and interactions with medications. Incorporating ginseng into a healthy lifestyle under professional guidance can harness its therapeutic potential effectively. Whether seeking to boost vitality, support immune function, or enhance mental clarity, ginseng stands as a testament to the enduring synergy between nature and health.
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|
__label__pos
| 0.981612 |
Climatology Analysis of the Daytime Topside Ionospheric Diffusive O+ Flux Based on Incoherent Scatter Radar Observations at Millstone Hill
Yihui Cai, Wenbin Wang, Shun Rong Zhang, Xinan Yue, Zhipeng Ren, Huixin Liu
研究成果: ジャーナルへの寄稿学術誌査読
6 被引用数 (Scopus)
抄録
This paper reports the characteristics of the topside ionospheric O+ diffusive flux ((Formula presented.)) during both geomagnetically quiet (0 ≤ Kp ≤ 2) and moderate (2 < Kp ≤ 4) times using incoherent scatter radar observations at Millstone Hill (42.6°N, 288.5°E) for solar minimum from 1970 to 2018. (Formula presented.) partially characterizes plasma mass exchange between the upper and lower part of the topside ionosphere through diffusion and sometimes serves as upper boundary conditions for ionosphere-thermosphere models. The altitude where the flux sign changes (mainly during daytime) is termed the transition height and the time when the flux sign changes (mainly at dawn and dusk) is termed the transition time. At quiet times, the daytime transition height is ∼100 km above the (Formula presented.) peak height (hmF2) in summer, and it is about 50 km above hmF2 in other seasons; the transition time is before 18 solar local time (SLT) in spring and winter, but after 18 SLT in summer and autumn. The daytime average upward (Formula presented.) above the transition height shows a significant seasonal variation with a minimum of (Formula presented.) in summer and a maximum of (Formula presented.) in autumn. Under geomagnetically moderate conditions, the transition height increases by ∼20 km in spring, winter, and autumn, but moves up by about 20–50 km in summer. The transition time occurs later by ∼1 hr in summer but ∼1 hr earlier in other seasons. The mean upward (Formula presented.) peaks in summer and minimizes in spring.
本文言語英語
論文番号e2021JA029222
ジャーナルJournal of Geophysical Research: Space Physics
126
10
DOI
出版ステータス出版済み - 10月 2021
!!!All Science Journal Classification (ASJC) codes
• 宇宙惑星科学
• 地球物理学
フィンガープリント
「Climatology Analysis of the Daytime Topside Ionospheric Diffusive O+ Flux Based on Incoherent Scatter Radar Observations at Millstone Hill」の研究トピックを掘り下げます。これらがまとまってユニークなフィンガープリントを構成します。
引用スタイル
|
__label__pos
| 0.591724 |
text: golang.org/x/text/message/catalog Index | Files
package catalog
import "golang.org/x/text/message/catalog"
Package catalog defines collections of translated format strings.
This package mostly defines types for populating catalogs with messages. The catmsg package contains further definitions for creating custom message and dictionary types as well as packages that use Catalogs.
Package catalog defines various interfaces: Dictionary, Loader, and Message. A Dictionary maintains a set of translations of format strings for a single language. The Loader interface defines a source of dictionaries. A translation of a format string is represented by a Message.
Catalogs
A Catalog defines a programmatic interface for setting message translations. It maintains a set of per-language dictionaries with translations for a set of keys. For message translation to function properly, a translation should be defined for each key for each supported language. A dictionary may be underspecified, though, if there is a parent language that already defines the key. For example, a Dictionary for "en-GB" could leave out entries that are identical to those in a dictionary for "en".
Messages
A Message is a format string which varies on the value of substitution variables. For instance, to indicate the number of results one could want "no results" if there are none, "1 result" if there is 1, and "%d results" for any other number. Catalog is agnostic to the kind of format strings that are used: for instance, messages can follow either the printf-style substitution from package fmt or use templates.
A Message does not substitute arguments in the format string. This job is reserved for packages that render strings, such as message, that use Catalogs to selected string. This separation of concerns allows Catalog to be used to store any kind of formatting strings.
Selecting messages based on linguistic features of substitution arguments
Messages may vary based on any linguistic features of the argument values. The most common one is plural form, but others exist.
Selection messages are provided in packages that provide support for a specific linguistic feature. The following snippet uses plural.Select:
catalog.Set(language.English, "You are %d minute(s) late.",
plural.Select(1,
"one", "You are 1 minute late.",
"other", "You are %d minutes late."))
In this example, a message is stored in the Catalog where one of two messages is selected based on the first argument, a number. The first message is selected if the argument is singular (identified by the selector "one") and the second message is selected in all other cases. The selectors are defined by the plural rules defined in CLDR. The selector "other" is special and will always match. Each language always defines one of the linguistic categories to be "other." For English, singular is "one" and plural is "other".
Selects can be nested. This allows selecting sentences based on features of multiple arguments or multiple linguistic properties of a single argument.
String interpolation
There is often a lot of commonality between the possible variants of a message. For instance, in the example above the word "minute" varies based on the plural catogory of the argument, but the rest of the sentence is identical. Using interpolation the above message can be rewritten as:
catalog.Set(language.English, "You are %d minute(s) late.",
catalog.Var("minutes",
plural.Select(1, "one", "minute", "other", "minutes")),
catalog.String("You are %[1]d ${minutes} late."))
Var is defined to return the variable name if the message does not yield a match. This allows us to further simplify this snippet to
catalog.Set(language.English, "You are %d minute(s) late.",
catalog.Var("minutes", plural.Select(1, "one", "minute")),
catalog.String("You are %d ${minutes} late."))
Overall this is still only a minor improvement, but things can get a lot more unwieldy if more than one linguistic feature is used to determine a message variant. Consider the following example:
// argument 1: list of hosts, argument 2: list of guests
catalog.Set(language.English, "%[1]v invite(s) %[2]v to their party.",
catalog.Var("their",
plural.Select(1,
"one", gender.Select(1, "female", "her", "other", "his"))),
catalog.Var("invites", plural.Select(1, "one", "invite"))
catalog.String("%[1]v ${invites} %[2]v to ${their} party.")),
Without variable substitution, this would have to be written as
// argument 1: list of hosts, argument 2: list of guests
catalog.Set(language.English, "%[1]v invite(s) %[2]v to their party.",
plural.Select(1,
"one", gender.Select(1,
"female", "%[1]v invites %[2]v to her party."
"other", "%[1]v invites %[2]v to his party."),
"other", "%[1]v invites %[2]v to their party.")
Not necessarily shorter, but using variables there is less duplication and the messages are more maintenance friendly. Moreover, languages may have up to six plural forms. This makes the use of variables more welcome.
Different messages using the same inflections can reuse variables by moving them to macros. Using macros we can rewrite the message as:
// argument 1: list of hosts, argument 2: list of guests
catalog.SetString(language.English, "%[1]v invite(s) %[2]v to their party.",
"%[1]v ${invites(1)} %[2]v to ${their(1)} party.")
Where the following macros were defined separately.
catalog.SetMacro(language.English, "invites", plural.Select(1, "one", "invite"))
catalog.SetMacro(language.English, "their", plural.Select(1,
"one", gender.Select(1, "female", "her", "other", "his"))),
Placeholders use parentheses and the arguments to invoke a macro.
Looking up messages
Message lookup using Catalogs is typically only done by specialized packages and is not something the user should be concerned with. For instance, to express the tardiness of a user using the related message we defined earlier, the user may use the package message like so:
p := message.NewPrinter(language.English)
p.Printf("You are %d minute(s) late.", 5)
Which would print:
You are 5 minutes late.
This package is UNDER CONSTRUCTION and its API may change.
Index
Package Files
catalog.go dict.go
Variables
var ErrNotFound = errors.New("catalog: message not found")
ErrNotFound indicates there was no message for the given key.
type Catalog Uses
type Catalog struct {
// contains filtered or unexported fields
}
A Catalog holds translations for messages for supported languages.
func New Uses
func New(opts ...Option) *Catalog
New returns a new Catalog.
func (*Catalog) Context Uses
func (c *Catalog) Context(tag language.Tag, r catmsg.Renderer) *Context
Context returns a Context for formatting messages. Only one Message may be formatted per context at any given time.
func (*Catalog) Languages Uses
func (c *Catalog) Languages() []language.Tag
Languages returns all languages for which the Catalog contains variants.
func (*Catalog) Set Uses
func (c *Catalog) Set(tag language.Tag, key string, msg ...Message) error
Set sets the translation for the given language and key.
When evaluation this message, the first Message in the sequence to msgs to evaluate to a string will be the message returned.
func (*Catalog) SetMacro Uses
func (c *Catalog) SetMacro(tag language.Tag, name string, msg ...Message) error
SetMacro defines a Message that may be substituted in another message. The arguments to a macro Message are passed as arguments in the placeholder the form "${foo(arg1, arg2)}".
func (*Catalog) SetString Uses
func (c *Catalog) SetString(tag language.Tag, key string, msg string) error
SetString is shorthand for Set(tag, key, String(msg)).
type Context Uses
type Context struct {
// contains filtered or unexported fields
}
A Context is used for evaluating Messages. Only one Message may be formatted per context at any given time.
func (*Context) Execute Uses
func (c *Context) Execute(key string) error
Execute looks up and executes the message with the given key. It returns ErrNotFound if no message could be found in the index.
type Message Uses
type Message interface {
catmsg.Message
}
A Message holds a collection of translations for the same phrase that may vary based on the values of substitution arguments.
func String Uses
func String(name string) Message
String specifies a plain message string. It can be used as fallback if no other strings match or as a simple standalone message.
It is an error to pass more than one String in a message sequence.
func Var Uses
func Var(name string, msg ...Message) Message
Var sets a variable that may be substituted in formatting patterns using named substitution of the form "${name}". The name argument is used as a fallback if the statements do not produce a match. The statement sequence may not contain any Var calls.
The name passed to a Var must be unique within message sequence.
type Option Uses
type Option func(*options)
An Option configures Catalog behavior.
Package catalog imports 6 packages (graph) and is imported by 3 packages. Updated 2017-10-16. Refresh now. Tools for package owners.
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ROS 2 Open Source Packages
Note
To use any of these packages, you must already have ROS 2 and the X-Series Arm packages installed. If you do not have these installed, follow the steps detailed in the ROS 2 Interface Setup.
Below is a list of all ROS 2 packages meant to be used with the X-Series robotic arms sold by Trossen Robotics. Packages were tested on Ubuntu Linux 20.04 and 22.04 using ROS 2 Galactic and Humble & Rolling respectively. Additionally, all ROS 2 nodes were written using Python or C++. However, any programming language capable of sending ROS 2 messages can be used to control the robots. The core packages inside this repo are as follows:
There are also several packages demonstrating possible applications of the core packages. A list of those packages is below in order of importance:
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@article {10.7554/eLife.20930, article_type = {journal}, title = {γ-Protocadherin structural diversity and functional implications}, author = {Goodman, Kerry Marie and Rubinstein, Rotem and Thu, Chan Aye and Mannepalli, Seetha and Bahna, Fabiana and Ahlsén, Göran and Rittenhouse, Chelsea and Maniatis, Tom and Honig, Barry and Shapiro, Lawrence}, editor = {Weis, William I}, volume = 5, year = 2016, month = {oct}, pub_date = {2016-10-26}, pages = {e20930}, citation = {eLife 2016;5:e20930}, doi = {10.7554/eLife.20930}, url = {https://doi.org/10.7554/eLife.20930}, abstract = {Stochastic cell-surface expression of α-, β-, and γ-clustered protocadherins (Pcdhs) provides vertebrate neurons with single-cell identities that underlie neuronal self-recognition. Here we report crystal structures of ectodomain fragments comprising cell-cell recognition regions of mouse γ-Pcdhs γA1, γA8, γB2, and γB7 revealing \textit{trans}-homodimers, and of C-terminal ectodomain fragments from γ-Pcdhs γA4 and γB2, which depict \textit{cis}-interacting regions in monomeric form. Together these structures span the entire γ-Pcdh ectodomain. The \textit{trans}-dimer structures reveal determinants of γ-Pcdh isoform-specific homophilic recognition. We identified and structurally mapped \textit{cis}-dimerization mutations to the C-terminal ectodomain structures. Biophysical studies showed that Pcdh ectodomains from γB-subfamily isoforms formed \textit{cis} dimers, whereas γA isoforms did not, but both γA and γB isoforms could interact in \textit{cis} with α-Pcdhs. Together, these data show how interaction specificity is distributed over all domains of the γ-Pcdh \textit{trans} interface, and suggest that subfamily- or isoform-specific \textit{cis}-interactions may play a role in the Pcdh-mediated neuronal self-recognition code.}, keywords = {clustered protocadherins, crystal structure, protein-protein recognition, protein diversity, neuronal self-avoidance}, journal = {eLife}, issn = {2050-084X}, publisher = {eLife Sciences Publications, Ltd}, }
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简单介绍：实现小程序授权登录功能
来自：互联网
时间：2022-11-07
阅读：
免费资源网 - https://freexyz.cn/
在我们平时工作、学习、生活中，微信小程序已成为我们密不可分的一部分，我们仔细留意下，每当我们使用一个新的小程序时，总会遇到如下页面:
简单介绍：实现小程序授权登录功能
这便是微信小程序授权登录功能了，授权登录后，我们就可以正常使用小程序，而小程序也会获取到我们的用户权益，手机号等个人信息
授权登录功能剖析
微信小程序的授权登录具体步骤如下所示
简单介绍：实现小程序授权登录功能
具体实现主要有以下三个步骤:
• 调用wx.login() 微信api获取临时登录凭证code，并回传到开发者服务器
• 调用auth.code2Session 微信api接口，获取用户唯一标识OpenID、 用户在微信开放平台帐号下的唯一标识UnionID和会话密钥session_key
• 通过步骤2获取的参数进行解密操作，获取用户手机号，头像等特性，并把需要的数据保存到缓存中
步骤实现代码如下：
一、获取临时登录凭证code
由于微信官方修改了getUserInfo接口，现在无法实现一进入微信小程序就会自动弹出授权窗口，所以我们只能通过button按钮让用户手动触发
我们先写一个简单的弹框，用isShow变量控制，isShow取决于步骤3中的缓存信息，当所有步骤都走通，会正确缓存用户信息，此时弹框隐藏，否则弹框都为显示状态
简单介绍：实现小程序授权登录功能
<view class='wx_dialog' wx:if="{{isShow}}">
<view class='wx_content'>
<text>需要先授权获取个人信息</text>
<button class="btn" open-type="getPhoneNumber" type="primary" bindgetphonenumber="getUserInfo">微信账号快速授权</button>
</view>
</view>
点击按钮时，调用getUserInfo方法，isShow设置为false,同时使用wx.login获取到登录凭证code
getUserInfo:e=>{ this.setData({ isShow:false
})
wx.login({ success: function (res) { let code = res.code // 登录凭证code
}
})
}
二、根据登录凭证code，获取用户登录信息
拿到登录凭证code后，调用auth.code2Session 微信api接口（此处为服务端操作，后端大佬搞定，我们直接调用他给的接口就好）
wx.request({ url: 获取用户信息的auth.code2Session微信api接口, method: 'POST', data:{ code:code//登录凭证code
}, header: { 'content-type': 'application/json;charset=UTF-8'
},
success: function (res) { var userphone= res.data.data //解密手机号
var msg = e.detail.errMsg; var sessionKey = userphone.session_key;//会话密钥
var encryptedData=e.detail.encryptedData; //签名
var unionid = userphone.unionid//唯一标识
var iv= e.detail.iv; //授权成功
if (msg == 'getPhoneNumber:ok') {
wx.checkSession({ success:function(){ //进行请求服务端解密手机号
this.deciyption(sessionKey,encryptedData,iv,unionid);
}
})
}
}
})
}
})
此时大多数用户信息我们已经获取了，但用户手机号，用户头像等信息还处于加密状态，我们需要去解密获取这些参数
三、根据用户信息，解密获取用户手机号
deciyption(sessionKey,encryptedData,iv,unionid){ var that = this;
wx.request({ url: 解密接口, method: 'POST', data: { sessionKey: sessionKey, encryptedData:encryptedData, iv: iv
}, header: { 'content-type': 'application/json;charset=UTF-8'
}, success: function(res) {
let data = res.data if (data.resultCode == 'success') {
wx.setStorageSync('userTel', data.data.phoneNumber);//存储解密后的用户手机号
}else{
wx.showToast({ title: '获取信息失败请重新授权', icon: 'none'
})
that.setData({ isShow:true
})
}
},
fail:function(res) {
wx.showToast({ title: '获取失败请重新授权', icon: 'none'
})
that.setData({ isShow:true
})
}
})
},
此时授权登录功能已完成
简单介绍：实现小程序授权登录功能
免费资源网 - https://freexyz.cn/
返回顶部
顶部
|
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开发者社区> 长风呼啸> 正文
阿里云
为了无法计算的价值
打开APP
阿里云APP内打开
如何生成安全的资源ID
简介：
+关注继续查看
一、前言
对于一些图片，文章，或者用户主页等，需要构造URL提供给外部。
对外发布URL时，通常是 “域名/路径/资源ID”,
其中，路径是可选项，比如生成短链接时可能就是直接“域名/资源ID”。
举例：
掘金的URL格式 ：
https://juejin.im/user/59abef0af265da246c4a3de1
https://juejin.im/post/5d8ab56df265da5bb252d67c
简书的 URL格式：
https://www.jianshu.com/u/11d3f06afbcd
https://www.jianshu.com/p/3df395d8a6bc
这些链接最后的资源ID部分是怎么构造的呢？
虽然无法确切知晓，但猜测一下也无妨。
掘金的资源ID，六进制字编码，32字节，可能时UUID（去掉分隔线）或者MD5
无论是UUID还是MD5，都是有随机性的，所以不担心被找规律，由于取值范围有128bit, 发生冲突的概率也微乎其微。
简书的资源ID，十六进制编码，12字节，也就是48bit，取值范围两百多万亿，够分配了，可读性也较好。
但是48bit的取值范围，就不然像UUID一样取随机值了，否则容易冲突，
因此，猜测原始ID是通过分段ID或者自增ID构造，自增ID的概率比较大，因为像snowflake算法这种需要的较长的有效位(bit数）。
但是自增ID的话，会呈现规律性（前面的字符不变），其他人就可以连续做几个请求，算出其自增的step，然后就可以穷举其所有资源了。
所以，在获取增长ID之后，十六进制编码之前，可能有一次加密。
那么，如何做加密呢？
以下讨论就不说“猜测”了，而是如果让我来做，我会如何设计。
二、对ID做加密
加密结果需要满足以下几点：
1、难以推测原ID；
2、加密后的ID和原ID一一映射，以避免重复；
3、对整数ID(64bit, 有效位<=64) 加密，加密后长度不变。
• 第1点：
MD5，SHA，AES，DES等都可以；
• 第2点：
MD5等摘要算法就不满足了，不同的原文可能计算出相同的哈希；
要满足一一映射，需要函数可逆。
• 第3点：
AES，DES也不满足。
AES加密结果最短也有16字节。
DES对小于8字节的原文加密，密文为8字节；对8字节的原文加密，密文为16字节。
综上，我们需要一个对整数的加解密，且密文不会“膨胀”的函数。
不妨去看下AES是怎么做的，抄个作业。
三、 AES 算法
接下来以AES128为例，看下AES的加密步骤。
部分图文来自文章 《码算法详解——AES》《AES算法描述及C语言实现》，侵删。
3.1 整体流程
AES128的核心运算是对16字节(128bit)的内容加解密。
若原文长于16字节，则进行分组，16字节一组，如果最后一组不足16字节，则补齐到16字节。
AES128要经历10轮运算，其中1-9轮是相同的，第10轮稍有一点区别。
每轮运算涉及字节替代、行移位、列混淆、轮密钥加等四个子运算；
每个子运算都有逆运算，用相反顺序逆运算可解密出原文。
3.2 字节替代（SubBytes）
字节替代需要用到S盒，S盒有两个数组，且命名为：SBOX， INV_SBOX。
S盒有这样的特点:
若y = SBOX[x]，则x = INV_SBOX[y]；
换个写法： x = INV_SBOX[SBOX[x]]。
举个栗子, 一个简单的2阶S盒：
int[] s_box = new int[]{2, 0, 3, 1};
int[] inv_s_box = new int[]{1, 3, 0, 2};
for (int x = 0; x < s_box.length; x++) {
System.out.print(inv_s_box[s_box[x]] + " ");
}
System.out.println();
for (int x = 0; x < s_box.length; x++) {
System.out.print(s_box[inv_s_box[x]] + " ");
}
输出：
0 1 2 3
0 1 2 3
AES的S盒是8阶的（8 x 8), 刚好取尽一个字节的范围(0-255)，至于怎么构造S盒本文就不展开了。
字节替代运算：
for (int j = 0; j < 16; j++) {
state[j] = SBOX[state[j]];
}
逆向字节替代:
for (int j = 0; j < 16; j++) {
state[j] = INV_S_BOX[state[j]];
}
字节替代运算提供了算法的混淆性
和代码混淆一样，原文本来具备可读性，混淆后就变得不可读了；
但是混淆后模式没有改变，比如原来方法foo()混淆后变为a(), 则混淆后所有调用foo()的地方都是a()。
也就是，字节替代后仍然存在统计特征。
3.3 行位移（ShiftRows）
行位移比较简单，就是将16字节作为一个4行4列的矩阵，
其中1、2、3行分别左移（逆运算为右移）1、2、3个字节。
image.png
3.4 列混淆（MixColumns）
列混淆是所有子运算中最复杂的。
和行位移一样，将16字节划分为4行4列；
不同的是，列混淆是分别对每一列的4个字节做运算（和一个4x4的矩阵左乘）。
解密时也是左乘矩阵，解密矩阵是加密矩阵的逆矩阵。
某一列的运算代码如下：
static inline uint8_t mul2(uint8_t a) {
return (a & 0x80) ? ((a << 1) ^ 0x1b) : (a << 1);
}
/*
* 左乘置换矩阵
* [d0] [02 03 01 01] [b0]
* [d1] = [01 02 03 01] . [b1]
* [d2] [01 01 02 03] [b2]
* [d3] [03 01 01 02] [b3]
*/
void multiply(uint8_t *b, uint8_t *d) {
uint8_t t = b[0] ^ b[1] ^ b[2] ^ b[3];
// d0 = (b0 << 1) + (b1 << 1) + b1 + b2 + b3 = ((b0 + b1) << 1) - b0 + t
d[0] = mul2(b[0] ^ b[1]) ^ b[0] ^ t;
d[1] = mul2(b[1] ^ b[2]) ^ b[1] ^ t;
d[2] = mul2(b[2] ^ b[3]) ^ b[2] ^ t;
d[3] = mul2(b[3] ^ b[0]) ^ b[3] ^ t;
}
矩阵运算需要进行加法运算和乘法运算，计算机的直接整数相加和直接整数相乘可能会溢出，从而丢失信息，不可逆；
所以AES引入了“伽罗华域（也叫伽罗瓦域）”运算，感兴趣的读者可以阅读文章《伽罗华域运算及C语言实现》了解一下。
行位移和列混淆共同提供了算法的扩散性
扩散的目的是让明文中的单个数字影响密文中的多个数字，从而使明文的统计特征在密文中消失。
最理想的情况是达到 雪崩效应 的效果:
当输入发生最微小的改变（例如，反转一个二进制位）时，也会导致输出的剧变（如，输出中一半的二进制位发生反转）。
如果只有列混淆运算，则最终的效果是 [0,3] , [4,7], [8,11], [12,15] 四个分组分别扩展；
加上了行位移，才可以达到[0, 15]的字节的扩散效果（明文一个字节改变，密文16个字节全都会随机变化）。
当然，需要经历多轮运算才会有雪崩效应的效果，只做一轮运算是不够的。
3.5 轮密钥加（AddRoundKey）
轮密钥加是四个字运算中最简单的，具体就是16个字节分别和轮密钥做异或运算
轮密钥是通过原始密钥计算得来的，AES128一共要做11次轮密钥加。
结合密钥的运算，提供了算法的保密性。
如果没有密钥运算，就好比门没上锁，再坚固的防盗门也是形同虚设。
四、整数加密方案
AES是典型SP型密码, SP型密码网络结构非常清晰，S被称为混淆层（非线性层），主要起混淆作用，P被称为扩散层，主要起扩散作用。
在AES算法中，
1、使用S盒做字节替代操作来达到混淆效果；
2、通过在伽罗瓦域做矩阵运算（MixColumns），对分组中的4个字节进行扩散；
3、同时结合行位移来交错MixColumns的子分组，从而实现整个分组的扩散；
4、最后混入“密钥”，实现算法的保密性。
值得注意的是，
AES128的分组大小是16字节，可以作为4x4的矩阵；
AES256的分组大小是32字节，可以作为4x8的矩阵。
MixColumns运算，可以看作时用4x4的矩阵左乘4xN的矩阵。
对于一个8字节的整数，我们也可以将其看作4x2的矩阵来做MixColumns运算。
所以，我们可以仿照AES的结构和运算，对一个Long类型的数值做加密运算：
public class LongEncoder {
private static final int ROUND = 8;
private static final byte[] S_BOX = {
99, 124, 119, 123, -14, 107, 111, -59
......
};
private static final byte[] KEY = {
-14, 40, 52, -119, -126, -47, 74, 73,
......
};
private static byte mul2(byte a) {
return (byte) (((a & 0x80) != 0) ? ((a << 1) ^ 0x1b) : (a << 1));
}
public static void mix_column(byte[] s, int i) {
byte t = (byte) (s[i] ^ s[1 + i] ^ s[2 + i] ^ s[3 + i]);
byte b0 = (byte) (mul2((byte) (s[i] ^ s[1 + i])) ^ s[i] ^ t);
byte b1 = (byte) (mul2((byte) (s[1 + i] ^ s[2 + i])) ^ s[1 + i] ^ t);
byte b2 = (byte) (mul2((byte) (s[2 + i] ^ s[3 + i])) ^ s[2 + i] ^ t);
byte b3 = (byte) (mul2((byte) (s[3 + i] ^ s[i])) ^ s[3 + i] ^ t);
s[i] = b0;
s[1 + i] = b1;
s[2 + i] = b2;
s[3 + i] = b3;
}
private static void shift_rows(byte[] state) {
byte t1 = state[7];
byte t0 = state[6];
state[7] = state[5];
state[6] = state[4];
state[5] = state[3];
state[4] = state[2];
state[3] = state[1];
state[2] = state[0];
state[1] = t1;
state[0] = t0;
}
public static long encode64(long value) {
byte[] state = long2Bytes(value);
for (int i = 0; i < ROUND; i++) {
for (int j = 0; j < 8; j++) {
int m = ((i << 3) + j);
// AddRoundKey and SubBytes
state[j] = S_BOX[(state[j] ^ KEY[m]) & 0xFF];
}
shift_rows(state);
mix_column(state, 0);
mix_column(state, 4);
}
for (int j = 0; j < 8; j++) {
state[j] ^= KEY[(ROUND << 3) + j];
}
return bytes2Long(state);
}
......
}
事实上，输入输出的长度（字节）不一定非要是4的倍数。
比如，对于6个字节的输入，可以这么处理：
public static long encode48(long value) {
byte[] state = long2Bytes(value);
for (int i = 0; i < ROUND; i++) {
for (int j = 0; j < 6; j++) {
int m = ((i << 3) + j);
// AddRoundKey and SubBytes
state[j] = S_BOX[(state[j] ^ KEY[m]) & 0xFF];
}
// 对于48bit的输入而言，就不需要ShiftRows了
// 因为先后对[0,3], [2,5]进行MixColumns已经可以对整个输入扩散了
mix_column(state, 0);
mix_column(state, 2);
}
for (int j = 0; j < 6; j++) {
state[j] ^= KEY[(ROUND << 3) + j];
}
// 输出的Long，高位的两个字节没有变
// 所以如果输入时小于2^48的数值，则输出也是小于2^48的数组
return bytes2Long(state);
}
再将加密后的数值进行进制转换，即可输出字符串形式的ID：
private static final byte[] HEX_DIGITS = {
'0', '1', '2', '3', '4', '5', '6', '7', '8', '9',
'a', 'b', 'c', 'd', 'e', 'f'};
/**
* 小于2^48的long数值转十六进制字符串
* @param n long类型整数
* @return 12字节的字符串（十六进制）
*/
public static String long48ToHex(long n) {
if((n >>> 48) > 0){
throw new IllegalArgumentException(n + " is bigger than 2^48");
}
byte[] buf = new byte[12];
for (int i = 5; i >= 0; i--) {
int b = (int) n;
int index = i << 1;
buf[index] = HEX_DIGITS[(b >> 4) & 0xF];
buf[index + 1] = HEX_DIGITS[b & 0xF];
n = n >>> 8;
}
return new String(buf);
}
至此，我们便完成了对整型ID的加密和编码。
至于整型ID的生成，网上有很多“ID生成方案”，我们这里就不作展开了。
五、总结
本文从网站URL入手，对资源ID做了简要的分析，文章用了大量篇幅讲AES的算法原理，稍有喧宾夺主的意味。
但就好比做包子，馅比面皮投入多也是正常的。
好吃的馅不单可以用来做包子，也可用于其他用途。
比如前些时候笔者写的一篇文章： 《漫谈唯一设备ID》，文中提到，需要根据主键ID计算出另一个的ID(作为UDID), 返回给客户端。
用本文的方法即可满足需求。
完整代码已上传github,
地址：https://github.com/No89757/LongEncrypt
参考资料
[1] 分组密码
[2] 伽罗华域运算及C语言实现
[3] 码算法详解——AES
[4] AES算法描述及C语言实现
[5] 雪崩效应
[6] 漫谈唯一设备ID
版权声明：本文内容由阿里云实名注册用户自发贡献，版权归原作者所有，阿里云开发者社区不拥有其著作权，亦不承担相应法律责任。具体规则请查看《阿里云开发者社区用户服务协议》和《阿里云开发者社区知识产权保护指引》。如果您发现本社区中有涉嫌抄袭的内容，填写侵权投诉表单进行举报，一经查实，本社区将立刻删除涉嫌侵权内容。
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__label__pos
| 0.924673 |
OpenStax: Microbiology
OpenStax: Microbiology
OpenStax: Microbiology
Lead Author(s): Openstax Content
Source: OpenStax
Student Price: FREE
An OpenStax textbook that presents the core concepts of microbiology with a focus on applications for careers in allied health.Download EPUB
Introduction
Figure 1. Malaria is a disease caused by a eukaryotic parasite transmitted to humans by mosquitos. Micrographs (left and center) show a sporozoite life stage, trophozoites, and a schizont in a blood smear. On the right is depicted a primary defense against mosquito-borne illnesses like malaria—mosquito netting. (credit left: modification of work by Ute Frevert; credit middle: modification of work by Centers for Disease Control and Prevention; credit right: modification of work by Tjeerd Wiersma)
Although bacteria and viruses account for a large number of the infectious diseases that afflict humans, many serious illnesses are caused by eukaryotic organisms. One example is malaria, which is caused by Plasmodium, a eukaryotic organism transmitted through mosquito bites. Malaria is a major cause of morbidity (illness) and mortality (death) that threatens 3.4 billion people worldwide.1 In severe cases, organ failure and blood or metabolic abnormalities contribute to medical emergencies and sometimes death. Even after initial recovery, relapses may occur years later. In countries where malaria is endemic, the disease represents a major public health challenge that can place a tremendous strain on developing economies.
Worldwide, major efforts are underway to reduce malaria infections. Efforts include the distribution of insecticide-treated bed nets and the spraying of pesticides. Researchers are also making progress in their efforts to develop effective vaccines.2 The President’s Malaria Initiative, started in 2005, supports prevention and treatment. The Bill and Melinda Gates Foundation has a large initiative to eliminate malaria. Despite these efforts, malaria continues to cause long-term morbidity (such as intellectual disabilities in children) and mortality (especially in children younger than 5 years), so we still have far to go.
Footnotes
1 Centers for Disease Control and Prevention. “Impact of Malaria.” September 22, 2015. http://www.cdc.gov/malaria/malaria_worldwide/impact.html. Accessed January 18, 2016.
2 RTS, S Clinical Trials Partnership. “Efficacy and safety of RTS,S/AS01 malaria vaccine with or without a booster dose in infants and children in Africa: final results of a phase 3, individually randomised, controlled trial.” The Lancet 23 April 2015. DOI: http://dx.doi.org/10.1016/S0140-6736(15)60721-8.
|
__label__pos
| 0.752387 |
Modular Regression Verification for Reactive Systems
Begutachtete Veröffentlichung in Tagungsband
Autor(en):Alexander Weigl, Mattias Ulbrich und Daniel Lentzsch
In:9th International Symposium on Leveraging Applications of Formal Methods, Verification and Validation (ISoLA 2020)
Verleger:Springer
Reihe:Lecture Notes in Computer Science
Band:12477
Teil:II: Engineering Principles
Jahr:2020
Seiten:25-43
DOI:10.1007/978-3-030-61470-6_3
BibTeX
@Inproceedings{Weigl2021a,
booktitle = {9th International Symposium on
Leveraging Applications of Formal Methods,
Verification and Validation ({ISoLA} 2020)},
series = {Lecture Notes in Computer Science},
publisher = {Springer},
title = {Modular Regression Verification for Reactive Systems},
author = {Alexander Weigl and Mattias Ulbrich and Daniel Lentzsch},
year = {2020},
series = {Lecture Notes in Computer Science},
volume = {12477},
part = {II: Engineering Principles},
editor = {Tiziana Margaria and Bernhard Steffen},
doi = {10.1007/978-3-030-61470-6_3},
pages = {25--43},
month = oct,
eventdate = {2020-10-20/2020-10-30},
venue = {Rhodes, Greece}
}
|
__label__pos
| 0.997732 |
Why The Cat Both Alive & Dead In Schrodinger's Experiment?
Welcome Guys.
Ever wondered why the cat both alive and dead at the same time before opening the box in Schrodinger's Quantum Mechanics Experiment. I bet our commonsense argues with the statement why both alive and dead, why not may be dead or may be not.
Right, lets obey with our commonsense, clearly cat may be alive or may be dead but not both. This happens only in our macro world. This never works with the quantum realm. In the Quantum world, the things are simultaneous and probabilistic.
Where This All Started?
It's all started with the imagination of our renowned Austrian Physicist Erwin Schrodinger. He imagined a cat put in a sealed box along with a radioactive substance, detector and a flask of poison. If the substance starts decaying, the detector releases a hammer which breaks the poison flask which finally kill the cat. But no one knows if the cat died or alive till opening the box. It's 50% dead and 50% alive.
the schrodinger's cat experiment
The Schrodinger's Cat Experiment (source: buzzle.com)
The Copenhagen interpretation of quantum mechanics implies that after a while, the cat is simultaneously alive and dead. Yet, when one looks in the box, one sees the cat either alive or dead, not both alive and dead. This poses the question of when exactly quantum superposition ends and reality collapses into one possibility or the other.
The Multi Worlds: Schrodinger's Cat Experiment
The-multi-world-possibilities
The Multi-world possibilities (source: wikipedia)
One of the greatest secrets of physics we still do not completely understand, is what happens when the smallest things interact with the big things, that is when quantum mechanics meets our everyday world? Schrodinger's cat experiment illustrates the multi world theories. If the cat was dead in one world, it may be alive in another and this happens simultaneously and infinitely in different worlds as of our assumptions.
That's it guys. This is what I got for you today. I will continue to post interesting articles on quantum physics, Subscribe us and be the first one to read it. :)
|
__label__pos
| 0.962797 |
ESCAPE RATE OF SURFACE-STATE ELECTRONS ON LIQUID-HELIUM BELOW 1-K
K KONO*, K KAJITA, S KOBAYASHI, W SASAKI
*Corresponding author for this work
Research output: Contribution to journalArticlepeer-review
7 Scopus citations
Abstract
The escape rate of surface state electrons on liquid helium was measured in the temperature interval 0.56<T<1.0 K for electron areal density below 5×107 cm−2, where the rate is independent of density. It is found that the escape rate is dominated by the scattering from the helium gas atoms in the vapor phase above 0.8 K, while it is dominated by that from ripplons (surface wave quanta) below 0.8 K. Quantitative agreement is obtained between our experiment and the theoretical calculation by Nagano et al.
Original languageEnglish
Pages (from-to)195-203
Number of pages9
JournalJournal of Low Temperature Physics
Volume46
Issue number3-4
DOIs
StatePublished - Feb 1982
Fingerprint
Dive into the research topics of 'ESCAPE RATE OF SURFACE-STATE ELECTRONS ON LIQUID-HELIUM BELOW 1-K'. Together they form a unique fingerprint.
Cite this
|
__label__pos
| 0.854421 |
DeepGrooveBallBearing
Model ElementSpecifies an analytical bearing that overtakes the radial and axial loads in both directions.
Class Name
DeepGrooveBallBearing
Description
The bearing is defined by the following dimensions:
Figure 1.
Attributes
The table below summarizes the attributes of the DeepGrooveBallBearing classes.
Variable Type Notes Figure Symbol
id Int() Unique identification number.
label Str() Descriptor of the bearing element.
rm Reference("Marker", required=True) The reference marker of the bearing element.
This marker points to the center of the bearing.
Defines the orientation of the bearing (z axis of reference marker is the bearing axial axis). The rm part must always be on the housing of the system.
number_of_rollers Int(required=True) The number of rolling elements in the bearing.
pitch_diameter Double(required=True) The pitch diameter of the bearing. dp
width Double(required=True) The width of the bearing. b
inner_diameter Double(required=True) The inner diameter of the bearing. d
inner_shoulder_diameter Double() The inner ring shoulder diameter. The default value is calculated internally as a function of pitch and roller diameters. d1
outer_diameter Double(required=True) The outer diameter of the bearing. D
outer_shoulder_diameter Double() The outer ring shoulder diameter. The default value is calculated internally as a function of pitch and roller diameters. D2
roller_diameter Double(required=True) The rolling element diameter. dr
bearing_clearance Double(0.0) The bearing clearance. Apply the correct clearance for each bearing.
inner_race_conformity Double(0.52) This value expresses the relation of the rolling element diameter to the inner raceway groove radius:
(fi: inner race conformity fi = ri/dr, ri: inner raceway radius, 0.51 < fi < 0.54 )
outer_race_conformity Double(0.53) This value expresses the relation of the rolling element diameter to the outer raceway groove radius:
(fo: outer race conformity fo = ro/dr, ro: outer raceway radius, 0.51 < fo < 0.54)
young_modulus Double() The young modulus of the bearing components. See Comment 2.
poisson_ratio Double() The Poisson ratio of the bearing components. See Comment 2.
bearing_density Double() The density of the bearing parts. See Comment 2.
roller_pass_frequency Bool (True) Defines if the bearing rotates with the inner and outer shaft. Default is set to True. Deactivating the rotation improves the computational performance, but neglects excitations caused by rollers passing through the loading zone.
friction_torque Bool (False) Activates the friction torque in the bearing.
simplified_graphics Bool (False) Creates simplified graphics for bearing rings and avoids graphics for rollers. This can reduce the H3D file size for large models.
rollers_force_graphics Bool (False) Activates force vector graphics in the rolling elements.
vec_gra_scale Double(1.0) Scales the rollers force vector graphics.
damping_force Bool (True) Activates the damping force in the bearing.
output_rm Reference("Marker") Defines the reference marker where the output results are written.
The default value is the defined bearing reference marker (rm).
damping_ratio Double(0.1) The damping ratio of the bearings. This value must be between 0.0-1.0.
static_load_rating Double(0.0) The static load rating of the bearing. It is used in the friction torque calculation.
lubricant_viscosity Double(0.0) The lubricant viscosity of the bearing in [cSt]. It is used in the friction torque calculation.
lubrication_method Enum("grease
oil_mist
oil_bath
oil_jet",
default="grease")
The lubrication method of the bearing, used in the friction torque calculation.
transition_velocity Double(0.1) The angular velocity at which the friction torque takes full effect [rad/sec]. Friction is ramped up smoothly from 0 until the transition velocity for solver stability.
inner_connection_part Reference("Part", required=True) The part/shaft that connects to the inner ring. See Comment 1.
outer_connection_part Reference("Part", required=True) The part /shaft that connects to the outer ring. See Comment 1.
na Int(60) Defines the number of nodes in the perimeter of ring graphics.
nc Int(6) Defines the number of nodes in the raceway of ring graphics.
no Int(2) Defines the number of nodes in the width of ring graphics.
Example
b = DeepGrooveBallBearing(
id=1,
label='bearing',
rm=Marker(part=ground),
number_of_rollers=10,
pitch_diameter=36,
width=12.0,
roller_diameter=6.746,
inner_diameter=25.0,
inner_shoulder_diameter=32.04,
outer_diameter=47.0,
outer_shoulder_diameter=39.96,
bearing_clearance=10e-3,
inner_race_conformity=0.52,
outer_race_conformity=0.53,
young_modulus=210000,
poisson_ratio=0.3,
bearing_density= 7.85e-06,
roller_pass_frequency=True,
friction_torque=True,
damping_force=True,
damping_ratio=0.1,
static_load_rating=6.55e3,
lubricant_viscosity=80,
lubrication_method='grease',
inner_connection_part=b_shaft,
outer_connection_part=B_Ground,
)
|
__label__pos
| 0.980638 |
and pdfWednesday, May 12, 2021 6:25:38 AM4
Newtonian And Non Newtonian Fluids Pdf
newtonian and non newtonian fluids pdf
File Name: newtonian and non newtonian fluids .zip
Size: 2180Kb
Published: 12.05.2021
Newton’s Law of Viscosity, Newtonian and Non-Newtonian Fluids
Instead of changing viscosity according to temperature, they thicken or thin in response to the amount of force applied to them.
The box accelerates at 1. A non-Newtonian fluid is a fluid that does not follow Newton's law of viscosity, i. In non-Newtonian fluids, viscosity can change when under force to either more liquid or more solid. The key difference between Newtonian and non Newtonian fluids is that Newtonian fluids have a constant viscosity, whereas non-Newtonian fluids have a variable viscosity.
A fluid that departs from the classic linear newtonian relation between stress and shear rate. I have recently been researching the difference between Newtonian and Non-Newtonian Fluids. Cream 5. Non —Newtonian fluids 2. A Newtonian fluid's viscosity remains constant, no matter the amount of shear applied for a constant temperature..
We gladly provide technical assistance to businesses in Wisconsin and upper Michigan. What is shear? If the viscosity of a fluid varies with the thermal or time history, the fluid is called non-Newtonian. Non-Newtonian fluids are polymers. Most fluids we know are non-Newtonian fluids. Sorry I'm a little late on posting. Fluids can be classified into two types depending on the viscosity as Newtonian fluids and non-Newtonian fluids. Though the differences between the three main states of matter seem to be clearly defined, there are some materials that require a closer look.
To take the vi… Non-Newtonian fluids are just the opposite — if enough force is applied to these fluids, their viscosity will change.
It is first important to explain some basic terms essential to the understanding of the difference between these two types of fluids. These fluids have a linear relationship between viscosity and shear stress. Most of the hydrocarbon fluids are Newtonian fluid, but in some cases, the fluid in the flowline should be considered to be non-Newtonian. First, what is the difference between Newtonian and non-Newtonian fluids? Paint The fluids that do … Thus, these fluids are a rare class of non-Newtonian fluids.
Usually, no real fluids fit the definition exactly. Non-Newtonian Fluids. A poly- mer is composed of long chains of repeated units known as monomers that are strung together to yield giant macromolecules. What are the differences between newtonian and non newtonian 1. Non-newtonian fluid. But we can assume some common liquids and gases such as water and air as Newtonian fluids. A comparison is made between these fluids and Newtonian fluids.
Also, they show an increased viscosity upon agitation. Terms of Use and Privacy Policy: Legal. Because the cornstarch solution's viscosity changes with an applied force, it is a non-Newtonian fluid.
You can probably guess that non-Newtonian fluids are the opposite of Newtonian fluids. Lubricating oil is more viscous than water but it is less dense, and mercury is very dense but has a low viscosity.
Non-Newtonian fluids such as ketchup change viscosity when force is applied. These are shear stress and shear rate. All fluids can be broken down into two basic types, Newtonian, and non-Newtonian. There are different behavioural characteristics of non-Newtonian fluids including viscoelasticity, time-dependent viscosity, etc.
Rheology concerns itself with how materials principally liquids, but also soft solids flow under applied forces. Fill out the form below to receive a PDF of the infographic, "Newtonian vs. Non-Newtonian Fluids" by email. Perhaps one would feel comfortable with this assumption for crude oil and refined products. Although most fluids we know are non-Newtonian fluids, water and air are considered as Newtonian fluids at normal conditions. Toothpaste 7. Ink 8.
Rheopectic fluids are liquids or gases whose viscosity of the fluid increases with stress over time. The viscosity of Newtonian fluids is a function of temperature and pressure.
Viscosity in newtonian and non newtonian fluids 1. Shear is the relative motion between adjacent layers of a moving fluid. The faster the ball falls, the lower the viscosity of the fluid. One way to visualize viscosity is by watching a metal ball fall through a glass container of the liquid at different speeds for various viscosities.
There are four different types of non-Newtonian fluids, in two main categories. Newtonian fluids have a linear relationship between the rate of deflection of a chunk of fluid and the applied shear stress.
There are many types of non-Newtonian fluids and this is a vast topic, which is studied under a branch of physics known as rheology. Coussot, P. Furthermore, when considering the shear rate and shear stress, in Newtonian fluids, we can observe a zero shear rate at zero shear stress.
Why do you need to know the difference? However, almost all salts, molten polymer material, blood, toothpaste, paint, corn starch and many other varieties of fluids are non-Newtonian fluids. If you're sizing or selecting pumps, mixers, or any other type of equipment that applies shear to fluid, these are all terms you should know. Example of Newtonian fluid: Air, water, Kerosene etc types of fluid.
This makes the behavior of non-Newtonian fluids very strange. One of the first things young science students learn is the difference between solids, liquids, and gases. However, most fluids we know have a variable viscosity.
Non-Newtonian fluids are the fluids that have a variable viscosity and a variable relationship with shear stress. Shear Stress: Shear stress is the force applied to an object fluid in this case divided by the area parallel to the force. The viscosity of these fluids can change under force.
On the assumption that pseudoplastic flow is associated with the formation and rupture of structural linkages a new flow equation is derived.
Instead of responding to temperature, the viscosity of a non-Newtonian fluid changes with the amount of pressure applied. Shear Stress: Shear stress is the… That means; the shear rate is directly proportional to shear stress. A key statement is "The viscosity is a function only of the condition of the fluid, particularly its temperature. Crane Engineering. In a strict sense, a fluid is any state of matter that is not a solid, and a solid is a state of matter that has a unique stress-free state.
Viscosity of time independent non-Newtonian fluids is a function of shear rate, but at the same time it has no memory of kinematic history. Yogurt 2. Compare the Difference Between Similar Terms.
It's important to fully understand the properties of the fluids you're transferring, mixing, or pumping because viscosity plays a major role in sizing and selecting equipment. For Newtonian fluids, viscosity is a constant of proportionality between shear rate and shear stress. Ask us about it! We can divide fluids, i. Non-Newtonian fluids include catsup, paint, liquid detergent, liquid polymers and a variety of other liquids. Newtonian fluid have constant viscosity and it does not rely upon the shear stress applied on it.
In other words, the ratio of the shear stress to the shear rate is constant throughout the fluid. Water is an example of a Newtonian fluid. Honey 3. However, non-Newtonian fluids are relatively common, and include oobleck which becomes stiffer when vigorously sheared , or non-drip paint which becomes thinner when sheared. A shear stress will not affect the viscosity of a Newtonian fluid. In Figure 5, there is a layer of fluid between a fixed lower plate and a moving upper plate.
These are aptly categorized as non-Newtonian fluids, and they have different properties than Newtonian fluids. Need help with a viscous situation? Viscosity is the state of being thick and sticky due to the internal friction of the fluid. Non-Newtonian fluids. A fluid is said to be Newtonian if its viscosity, which is the measure or ability of a fluid to resist flow, only varies as a response to changes in temperature or pressure.
A Newtonian fluid will take the shape of its container. These are summarized below: Unlike a Newtonian fluid, applying forces to this non-Newtonian fluid causes its particles to behave more like a solid. But does heavy crude at some point low temperature stop behaving as … The name Newtonian comes from Isaac Newton, who was the first scientist to use a differential equation to postulate the relationship between shear stress and shear rate of fluids.
The behaviour of these fluids can be described as a time-dependent dilatant behaviour. With Newtonian fluids, viscosity does not depend on the shear rate.
NON NEWTONIAN FLOW AND APPLIED RHEOLOGY
Skip to search form Skip to main content You are currently offline. Some features of the site may not work correctly. DOI: A non-Newtonian fluid is a fluid whose flow properties differ in many ways from those of Newtonian fluids. Most commonly the viscosity of non-Newtonian fluids is not independent of shear rate or shear rate history. In practice, many fluid materials exhibits non-Newtonian fluid behavior such as: salt solutions, molten, ketchup, custard, toothpaste, starch suspensions, paint, blood, and shampoo etc. View PDF.
A non-Newtonian fluid is a fluid that does not follow Newton's law of viscosity , i. In non-Newtonian fluids, viscosity can change when under force to either more liquid or more solid. Ketchup , for example, becomes runnier when shaken and is thus a non-Newtonian fluid. Many salt solutions and molten polymers are non-Newtonian fluids, as are many commonly found substances such as custard , [1] honey , [1] toothpaste , starch suspensions, corn starch , paint , blood , melted butter , and shampoo. Most commonly, the viscosity the gradual deformation by shear or tensile stresses of non-Newtonian fluids is dependent on shear rate or shear rate history.
Successful characterization of viscosity is key in determining if a fluid is Newtonian or non-Newtonian. Each unique material has its own behavior when subjected to flow, deformation or stress. Depending on their viscosity behavior as a function of shear rate, stress, deformation history Some examples of Newtonian fluids include water, organic solvents, and honey. For those fluids viscosity is only dependent on temperature.
Newton’s Law of Viscosity, Newtonian and Non-Newtonian Fluids
Instead of changing viscosity according to temperature, they thicken or thin in response to the amount of force applied to them.
Non-Newtonian fluid
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These two volumes contain chapters written by experts in such areas as bio and food rheology, polymer rheology, flow of suspensions, flow in porous media, electrorheological fluids, etc. Computational as well as analytical mathematical descriptions, involving appropriate constitutive equations deal with complex flow situations of industrial importance. This work is unique in that it brings together state of the art reviews and recent advances in a variety of areas, involving viscoelastic materials, in a desirable and timely manner. We are always looking for ways to improve customer experience on Elsevier.
Newtonian and Non-Newtonian Fluids: Velocity Profiles, Viscosity Data, and Laminar Flow Friction Factor Equations for Flow in a Circular Duct. January DOI.
Fill out the form below to receive a PDF of the infographic, "Newtonian vs. Non-Newtonian Fluids" by email. This colorful document identifies the differences between the two concepts in a brief, easy-to-understand format.
Viscosity of Newtonian and non-Newtonian Fluids
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Applied Rheology Volume 21 Issue 5 Abstract: The influence of short fibre addition on the rheological behaviour of different non-Newtonian fluids is investi-gated experimentally.
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This book provides an up-to-date overview of mathematical theories and research results in non-Newtonian fluid dynamics.
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| 0.63091 |
blob: 7d66d84264bb165f026b32efb25104d9d1759fc1 [file] [log] [blame]
/* Copyright (c) 2014, The Linux Foundation. All rights reserved.
*
* Redistribution and use in source and binary forms, with or without
* modification, are permitted provided that the following conditions are
* met:
* * Redistributions of source code must retain the above copyright
* notice, this list of conditions and the following disclaimer.
* * Redistributions in binary form must reproduce the above
* copyright notice, this list of conditions and the following
* disclaimer in the documentation and/or other materials provided
* with the distribution.
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* from this software without specific prior written permission.
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#ifndef _LOC_MISC_UTILS_H_
#define _LOC_MISC_UTILS_H_
#ifdef __cplusplus
extern "C" {
#endif
/*===========================================================================
FUNCTION loc_split_string
DESCRIPTION:
This function is used to split a delimiter separated string into
sub-strings. This function does not allocate new memory to store the split
strings. Instead, it places '\0' in places of delimiters and assings the
starting address of the substring within the raw string as the string address
The input raw_string no longer remains to be a collection of sub-strings
after this function is executed.
Please make a copy of the input string before calling this function if
necessary
PARAMETERS:
char *raw_string: is the original string with delimiter separated substrings
char **split_strings_ptr: is the arraw of pointers which will hold the addresses
of individual substrings
int max_num_substrings: is the maximum number of substrings that are expected
by the caller. The array of pointers in the above parameter
is usually this long
char delimiter: is the delimiter that separates the substrings. Examples: ' ', ';'
DEPENDENCIES
N/A
RETURN VALUE
int Number of split strings
SIDE EFFECTS
The input raw_string no longer remains a delimiter separated single string.
EXAMPLE
delimiter = ' ' //space
raw_string = "hello new user" //delimiter is space ' '
addresses = 0123456789abcd
split_strings_ptr[0] = &raw_string[0]; //split_strings_ptr[0] contains "hello"
split_strings_ptr[1] = &raw_string[6]; //split_strings_ptr[1] contains "new"
split_strings_ptr[2] = &raw_string[a]; //split_strings_ptr[2] contains "user"
===========================================================================*/
int loc_util_split_string(char *raw_string, char **split_strings_ptr, int max_num_substrings,
char delimiter);
/*===========================================================================
FUNCTION trim_space
DESCRIPTION
Removes leading and trailing spaces of the string
DEPENDENCIES
N/A
RETURN VALUE
None
SIDE EFFECTS
N/A
===========================================================================*/
void loc_util_trim_space(char *org_string);
#ifdef __cplusplus
}
#endif
#endif //_LOC_MISC_UTILS_H_
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| 0.906845 |
ICP-OES: Why spectral lines are true peaks and how this can fool the user.
ICP–OES is a common technique in analytical chemistry, which is characterized by simplicity, simultaneous multi-elemental determination capability, high sensitivity, linear dynamic range, low detection limits, and good precision [1]. It seems to be the perfect method for element analysis (esp. metals). However, sometimes it can fool the user.
ICP–OES instruments use an inductive–coupled plasma to excite atoms and ions, which subsequently relax and emit light [2]. The result is an atomic spectra of a sample. Atomic spectra are ‘line spectra’, because they consist – in contrast to the typical, continuous band spectra of molecules – of a set of isolated lines [3].
I think that the concept of ‘spectral lines’ is often misunderstood. The very name ‘spectral line’ is a bit misleading and may suggest an (nearly) infinitively thin line at a certain wavelength. This imagination is supported by both the quoted precision of listed wavelengths (3–4 digits, e.g. as in [4]) and the rare specification of the widths of these lines in publications.
It is true that these lines are very narrow (low picometer range). However, they are still peaks possessing a wideness and their profiles have not always a nice, symmetrical Gaussian shape [5]. Nearby peaks can overlap. This leads to spectral interferences [2]. If not aware of this fact the results may fool the user of ICP–OES.
For instant, high calcium concentrations in wine samples will increase the intensity for aluminium by a factor of 8 [6]. For visualization of this example, I plotted below the involved peaks for Al and Ca (all at around 396 nm), using the position and widths from Boumans et al. [5].
Figure 1. Plot showing two Ca lines and the Al line around 396 nm.
Figure 1. Plot showing two Ca lines and the Al line around 396 nm. Peak positions and widths were taken from Boumans et al.
At the first glance, these peaks are very narrow and not overlapping at all. But let us zoom in a little by scaling the y–axis. The result is shown below.
Figure 2. Plot showing two Ca lines and the Al line around 396 nm.
Figure 2. Plot showing a zoomed–in version of Figure 1 with two Ca lines and the Al line around 396 nm. Peak positions and widths were taken from Boumans et al.
As can be seen now, the bases of two peaks, viz. Ca II and Al , are clearly overlapping (also a bit of the other Ca peak is overlapping). The dashed line in this plot marks 396.152 nm, which is the typically line used for Al determination. In this plot (as before) the intensities of all lines are the same and the influence of both Ca lines on the Al line is negligible.
However, from the concentration used to determine line shapes [5] one can deduce that the intensities of the Ca lines are much higher than of the Al line. Relative to mass (0.5 µg/mL for Ca, 100 µg/mL for Al) the ratio is 200. Relative to amount of element (0.0125 µmol/mL for Ca, 3.7 µmol for Al) the ratio becomes almost 300.
Thus, I assumed that both Ca lines are 300–times more intense than the Al line for this example. Of course, this ratio would only apply if the same amount of element is present in the sample. Typical concentrations of both elements in wine are 70 mg/L (i.e. 1.750 mM) and about 0.9 mg/L (i.e. 0.033 mM) for Ca and Al, respectively [6]. Considering this values, the ratio becomes 15 750! I replotted the peaks using this ratio in the figure below.
Figure 3. Plot showing the two Ca lines and the Al line around 396 nm with adjusted intensity ratio of 15750.
Figure 3. Plot showing the two Ca lines and the Al line around 396 nm with adjusted intensity ratio of 15 750 (Intensities of Ca lines: 15 750 each. Intensity of Al line: 1). Blue circle marks the overlap of the Ca lines and the Al line. Peak positions and widths were taken from Boumans et al.
Now, the influence of the Ca lines on the Al line is not negligible anymore. In this case, both Ca lines increase the total signal of Al (at 396.153 nm) by about 48%. If further increasing the Ca concentration, at some point, the overall intensity at the Al line will mainly be defined by the Ca concentration. In the end, this will lead to the previously mentioned factor of 8 [6]. Note that the left Ca line has an non–negligible influence even it is about 3 nm away!
There are more examples for interfering elements, so be aware when using ICP–OES! It will (like every other spectroscopic method) fool every user, who is not aware of possible spectral overlaps. I myself was not aware of this problem since I coincidentally read about it! Now I also demonstrate and teach this to my students.
The good thing is: There are solutions! Of course, the solution depends on the actual analytical problem. One can choose alternative, interference–free lines, use multi–line detection schemes, or (if the source and its influence is fully known!) substract out the influences on the measured signal [2]. Also, spectral overlaps can be minimized (not eliminated!) by using high-resolution spectrometers [2].
References
1. K. Satyanarayana, and S. Durani, "Separation and inductively coupled plasma optical emission spectrometric (ICP-OES) determination of trace impurities in nuclear grade uranium oxide", Journal of Radioanalytical and Nuclear Chemistry, vol. 285, pp. 659-665, 2010. http://dx.doi.org/10.1007/s10967-010-0591-8
2. J.W. Olesik, "Elemental Analysis Using ICP-OES and ICP/MS", Analytical Chemistry, vol. 63, pp. 12A-21A, 1991. http://dx.doi.org/10.1021/ac00001a711
3. R.D. Cowan, "The theory of atomic structure and spectra", Open Library https://openlibrary.org/works/OL6439482W/The_theory_of_atomic_structure_and_spectra
4. doi:10.1039/C2JA30276D (Timed out)
5. doi:10.1016/0584-8547(86)80005-2 (Timed out)
6. doi:10.1007/BF01244684 (Timed out)
Kcite was unable to retrieve citation information for all the references, due to a timeout. This is done to prevent an excessive number of requests to the services providing this information. More references should appear on subsequent page views
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Optimization
Various optimization algorithms
Source Code
Source code can be found here
Preliminaries
• All tests are preformed with a sample size of 30
• Implementation of algorithms follows the canonical one found on Wikipedia , unless stated otherwise.
• Source Code has assertions that make sure every solution is valid
Traveling Salesman Problem
Problem Description and Motivation
Given a collection of nodes X={x1,,xn}X = \{x_1, \ldots, x_n\} with an associated set of weighted edges E={(xi,xj,di,j) xi,xjX}E = \{ (x_i, x_j, d_{i,j})| \space x_i, x_j \in X\} , where the weight di,jd_{i, j} is the distance between xix_i and xjx_j. Find the shortest path to visit each xix_i in XX and return to x1x_1 after. For our implementation we will also add the restriction that as max(di,j)=αmax(d_{i, j}) = \alpha, this makes density increase as we add more nodes instead of adding spacing.
An example of a real world scenario is suppose you run a logistics company wherein you have to ship morning packages to nn number of customers every morning. Of course, as a profit maximizing firm it is imperative for you to find the shortest distance to cut down on gas expenditure, billed driver hours, and to increase client satisfaction. One could modify the distance function to account for speed limits, traffic, etc easily allowing extensibility.
The first thought would be to just take XX and try every permutation of elements until you find the shortest path. However, this would result in having to check
n!=n(n1)1n! = n \cdot(n - 1) \ldots 1
number of solutions, while this may seem innocent 11!11! is approximately 4 billion, and hence unfeasible to calculate in this manner. Given the TSP problem can be classified as NP-Hard there is no deterministic algorithm that solves it in polynomial time - hence we typically rely on non-deterministic heuristics to find solutions after 10 nodes.
Why Use a Heuristic?
Suppose you see the problem starts to become intractable after you reach only 10 nodes and you decide to just randomly pick nodes to travel to instead. Below is a comparison between random picking vs the simplest heuristic out there: simply picking the closest node to travel to.
Above we see as the number of nodes increases randomly picking your next choice will lead to a highly unoptimal solution even compared to the simplest way of informed selection.
Possible Heuristics
• Branch and Bound
This is an exact algorithm that prunes paths if they go above the best distance found already. The problem is if many of the distances are close enough you will devolve into a brute force search, given you cannot know the chances of pruning upfront this makes this method unreliable as a standard heuristic - but it is strictly superior alternative to calculating the exact distance. One could improve the average running time by improving the pruning of the possible paths.
• Nearest Neighbor
The simplest possible (good) heuristic is continually picking the closest neighbor of our given node until we run out of options and then returning to our start node. This is a fast implementation that on average finds a solution only 25% worse than the optimal solution. Given how fast this algorithm is it is phenomenal at using as an initial guess for other heuristics. One could also try to optimize this algorithm by looking ahead and having some discount future function.
• Ant System
An interesting alternative is called the Ant System algorithm and it's derivative implementations. The Ant System is a learning algorithm that tries to mimic biological optimization used by ants. Given a collection of nodes XX the corresponding weighted edges also have a value which represents the amount "pheromones" present i.e. E={(xi,xj,(di,j,pi,j) xi,xjX}E = \{ (x_i, x_j, (d_{i,j}, p_{i, j})| \space x_i, x_j \in X\}. Each Ant deposits pheromones for each edge from the graph it selects wherein the better solutions deposit more pheromones on each edge. Each ants selects a node based on a function that takes into account both pheromone deposited and the distance of that given edge. H(p,d)H(p, d). In effect this means after a given number of iterations better solutions will be selected for and bad solutions will be pruned.
• Simulated Annealing
Simulated Annealing works by traversing state space to try and find the global minimum of some goal function, In our case we are looking to minimize the distance a tour will take. To traverse the state space one finds some sort of state transition for our current state and given an acceptance function we decide to accept this transition or not. This acceptance function is informed by a parameter called "temperature" which decays, such that we accept worse solutions less as time goes on. We always accept a better solution than our current one. For our implementation our state transition function is defined by selecting a random segment of our tour and then allowing two operations: reverse and splice. reverse simply reverses the segment in place while splice takes the segment and puts it into another portion of our path.
Comparison vs Exact
Above we see that Brute Force and it's improvement Branch and Bound start becoming intractable at 12 nodes compared to our heuristics, hence solving this problem exactly is not feasible. Also note that Simulated Annealing was able to find the exact solutions on each test.
Heuristic Comparison
Above we see that while the Ant System always performs better than the Neighbor in finding an optimal solution as density increases the time to compute the solution doubles roughly every 10 additional nodes.. While this is substantially better than the exact algorithm, Ant System still starts to become unfeasible as nn\rarr \infty. However, our implementation of Simulated Annealing increases linearly in time for small amounts of nodes which is much better, especially since Nearest Neighbor time to found a solution increases at roughly the same rate. Furthermore, Simulated Annealing performs better than the other two heuristics in minimizing error in finding a solution. In generality we can suppose that a learning heuristic like Ant System or Simulated Annealing could be useful as an alternative in an optimization problem.
Simulated Annealing at Scale
Using 50 to 500 Nodes
Above we see Simulated Annealing stays competitive as the number of nodes increases, It doesn't become much worse time wise to find a solution while staying strictly better at minimizing error while finding a solution. Hence we can see it would be a preferred heuristic to use, especially since it can be generalized to any optimization problem which you can boil down to traversing a state space where you can be informed by a loss/error function.
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Thermocapillary rupture in falling liquid films at moderate Reynolds numbers
Evgeny Shatskiy, Evgeny Chinnov
Research output: Contribution to journalArticlepeer-review
Abstract
An experimental study of the flow of a water film over a heated surface for Re = 15-50 was performed. The influence of the development of thermocapillary instability on the wave amplitudes, the deformation of the surface of the liquid film, and the formation of the first stable dry spot on the heater are investigated. It is shown that the interaction of waves with thermocapillary structures can lead to an increase in the critical heat flux corresponding to the rupture of the liquid film, as compared with the data known in the literature.
Original languageEnglish
Article number0055
JournalEPJ Web of Conferences
Volume159
DOIs
Publication statusPublished - 25 Oct 2017
Fingerprint Dive into the research topics of 'Thermocapillary rupture in falling liquid films at moderate Reynolds numbers'. Together they form a unique fingerprint.
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Buschhagen_Dissertation.pdf (48.57 MB)
Download file
Self-Excited Instability Regimes of Confined Turbulent Premixed Jet Flames
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thesis
posted on 04.03.2020, 13:23 by Timo BuschhagenTimo Buschhagen
The demand for more efficient and cleaner terrestrial gas turbine engines for energy generation has accelerated with stringent emission regulations. Most terrestrial gas turbine engine configurations utilize lean premixed jet stabilized flames for energy extraction. These highly turbulent flames provide the necessary thermal power densities and low NOx emissions but are prone to combustion instabilities. This work studies the stability of a canonical premixed turbulent jet flame to changes in operation condition at elevated pressure. The underlying coupling mechanisms leading to the excitation of different modes are of interest. Different fuel injection schemes are considered, to delineate the influence of system coupled excitation and flow field intrinsic mechanisms that lead to self-excitation of longitudinal and transverse modes in an axisymmetric jet flame.
The fundamental longitudinal mode sensitivity to operation conditions was investigated in a technically-premixed configuration, for which the fuel injector is exposed to the system dynamics. Global equivalence ratio fluctuations paired with flame-vortex interactions were observed to sustain the 1L mode, for which leaner operating conditions lead to an increase of the limit-cycle amplitude.
To solely focus on the shear layer dynamics involved in the feedback loop of the 1L instability a "fully-premixed" configuration (FPC) of the combustor was designed. The elimination of potential equivalence ratio fluctuations lead to the excitation of longitudinal as well as transverse, and spinning modes, which can be associated with specific burner operation envelope regions. 1L mode coupled flame dynamics indicate axisymmetric emission patterns in OH* emission imaging corresponding to axisymmetric instabilities in the shear layer. Transverse modes correlate with an asymmetric shear layer roll-up process and a flapping motion of the flame. Spinning modes are characterized by high levels of limit-cycle amplitude and a single wave is observed that travels in the annulus of the recirculation zone. From the high speed imaging an azimuthal wave speed of up to 90% of the Chapman-Jouguet velocity for the natural gas - air mixture is computed.
The transverse mode is found to be sensitive to changes in chamber pressure and injector velocity. For a baseline injector velocity at which 1T mode excitation occurs, an increase in system pressure lead to an increase of the 1T mode amplitude. The 1T mode excitation is found to be sensitive to the injector velocity, where the highest amplitude are observed for a Strouhal number band of 1.6-1.7 based on the injector diameter.
A linear stability analysis (LSA) of the underlying base flow field is performed in order to assess if the underlying shear layer instability modes determine the selection of the instability regime for a given flame condition. Two flow perturbation modes are supported by the flow field, Mode 1 is associated with the recirculation zone domain, where preferred mode frequencies favors coupling with longitudinal acoustics chamber modes. Mode 2 resides in the mixture jet, for which preferred mode frequencies match local 1T acoustic eigenfrequencies of the chamber. It is found that for 1T mode dominated operation cases, the recirculation zone associated mode is stabilized leading to to the excitation of the 1T mode in the injector near field. For operation cases showing predominantly longitudinal combustion instability, the recirculation zone mode shows elevated growth rates in the injector near field paired with preferred frequencies that are compatible with the longitudinal acoustic eigenmodes.
These observations indicate that the preferred shear layer perturbation modes impact the selection of the combustion instability regime for a given flame operation condition. These results can be utilized to build a model framework for the design of jet flame type burners, to avoid geometries and operation regions that show a high potential for transverse mode excitation by the underlying flow field.
History
Degree Type
Doctor of Philosophy
Department
Aeronautics and Astronautics
Campus location
West Lafayette
Advisor/Supervisor/Committee Chair
Dr. Carson D. Slabaugh
Additional Committee Member 2
Dr. William E. Anderson
Additional Committee Member 3
Dr. Rohan M. Gejji
Additional Committee Member 4
Dr. Robert P. Lucht
Additional Committee Member 5
Dr. Carlo Scalo
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What Is Rosacea? – (Symptoms & Causes)
0
875
rosacea
rosacea
Want the most effective way to get rid of acne long term - this stuff is AMAZING...
Rosacea occur when spoiled blood vessels affect facial skin. It is one kind of common disease that occurs due to disorder of the blood vessels. Among the total world population about 48 percent suffers from Rosacea. On the other hand, It is one of the most mistaken conditions of the skin. When you are affected by this, at that time you need to talk with persons who know more about it.
Acne Rosacea Symptoms
Change of the skin color into red is the common sign of rosacea. Sometimes the body portions become different and these are the nerves, forehead and nose. This kind of soreness and coloring of skin can spread to the scalp, ears, and chest. The pink shade can make a stable form as increases. Moreover there can be a clear sign of the small blood vessels at the skin surface, harsh, eyes and rotary brave that look like red. Most of these simple signs are fat noses. The first signs are irritating redness that is often incorrectly credited to cleaning, training and temperature differences. There are many people who are unclear about this through pimple vulgaris.
Rosacea Pictures
rosacea
Rosacea
What Causes Rosacea?
Usually it is noticed that the people with fair skin are the best weak section with mention to it. Thus, It becomes the cause of genetic strain. Descent people from the Celtic bonds are certainly liable to suffer from this redness like a pimple.
All ages of people can be affected by this. It is observed that people in the age of 30-50 can certainly be affected by rosacea. Women in their mid-ages are the most open part of the public. The main cause is the menopause helped of hot blooms. signs are simpler with situation to men.
We cannot find any accord among the medical agents as to the exact rosacea. There is an agreement in views of the people such as to the source of rosacea. It happens when limbs again and again increase the blood vessels and thus the blood vessels get hurt. Then the spoiled blood vessels dilate easily. Further they always remain opened for long time. The result is soreness of the hollow portion and its coloring.
There are also different situations that are causes of this and one of them are active actions of glowing and blushing. These glowing and blushing are the causes of cold weather, pressure, severe tan and great heat contact mainly from the sun. Moreover Rosacea can be begun by rapid changes in temperatures when traveling or in excited rooms mainly in winter.
Rosacea Papules
The papule can be the method of a boil, a pimple or a wave for that material. The intermediaries and secure cells emit from the skin bed which is really micro-vascular by nature and this occurs for rosacea. This leads to the stimulating pimple. Definite food items that contain high amount of histamine have been known as the blast of it’s signs among people. Hot food like alcoholic elements can certainly set off rosacea.
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Skip to main content
\(\require{cancel}\)
Physics LibreTexts
13.8: Improved Triangle Ratios
The equation of motion of the orbiting body is
\[\ddot{\textbf{r}} = - \frac{GM}{r^3} \textbf{r} . \label{13.8.1} \tag{13.8.1}\]
If we recall equation 13.4.2, this can be written
\[\ddot{\textbf{r}} = -k^2 \left( \frac{a^3}{r^3} \right) \textbf{r} . \label{13.8.2} \tag{13.8.2}\]
If we now agree to express \(r\) in units of \(a\) (i.e. in Astronomical Units of length (au)) and time in units of \(1/k\) \((1/k = 58.132 \ 440 \ 87\) mean solar days), this becomes merely
\[\ddot{\textbf{r}} = -\frac{1}{r^3} \textbf{r} . \label{13.8.3} \tag{13.8.3}\]
In these units, \(GM\) has the value \(1\).
Now write the \(x\)- and \(y\)- components of this equation, where \((x , \ y)\) are heliocentric coordinates in the plane of the orbit (see sections 13.5 or 10.7).
\[\ddot{x} = -\frac{x}{r^3} \label{13.8.4} \tag{13.8.4}\]
and \[\ddot{y} = - \frac{y}{r^3}, \label{13.8.5} \tag{13.8.5}\]
where \[ x^2 + y^2 = r^2 . \label{13.8.6} \tag{13.8.6}\]
The areal speed is \(\frac{1}{2}h = \frac{1}{2} \sqrt{GMl}\) or, in these units, \(\frac{1}{2} \sqrt{l}\) where l is the semi latus rectum of the orbit in \(\text{au}\)
Let the planet be at \((x, \ y)\) at time \(t\). Then at time \(t + δt\) it will be at \((x + δx , y + δy)\), where
\[δx = \dot{x} δt + \frac{1}{2!} \ddot{x} (δt)^2 + \frac{1}{3!} \dddot{x} (δt)^3 + \frac{1}{4!}\ddddot{x} (δt)^4 + ... \label{13.8.7} \tag{13.8.7}\]
and similarly for y.
Now, starting from equation \ref{13.8.4} we obtain
\[\dddot{x} = \frac{3x \dot{r}}{r^4} - \frac{\dot{x}}{r^3} \label{13.8.8} \tag{13.8.8}\]
and \[\ddot{\ddot{x}} = 3 \left( \frac{\dot{x}\dot{r}}{r^4} + \frac{x \ddot{r}}{r^4} - \frac{4x \dot{r}^2}{r^5} \right) - \frac{r^3 \ddot{x} - 3r^2 \dot{x} \dot{r}}{r^6} . \label{13.8.9} \tag{13.8.9}\]
(The comment in the paragraph preceding equation 3.4.16 may be of help here, in case this is heavy-going.)
Now \(\ddot{x}\) and \(x\) are related by equation \ref{13.8.4}, so that we can write equation \ref{13.8.9} with no time derivatives of \(x\) higher than the first, and indeed it is not difficult, because equation \ref{13.8.4} is just \(r^3 \ddot{x} = -x\). We obtain
and \[\ddot{\ddot{x}} = x \left( \frac{1}{r^6} - \frac{12 \dot{r}^2}{r^5} + \frac{3 \ddot{r}}{r^4} \right) + \frac{6 \dot{x} \dot{r}}{r^4} . \label{13.8.10} \tag{13.8.10}\]
In a similar fashion, because of the relation \ref{13.8.4}, all higher time derivatives of \(x\) can be written with no derivatives of \(x\) higher than the first, and a similar argument holds for \(y\).
Thus we can write equation \ref{13.8.7} as
\[x + δx = Fx + G \dot{x} \label{13.8.11} \tag{13.8.11}\]
and similarly for \(y\):
\[y + δy = Fy + G \dot{y} , \label{13.8.12} \tag{13.8.12}\]
where \[F = 1 - \frac{1}{2r^3} (δt)^2 + \frac{\dot{r}}{2r^4} (δt)^3 + \frac{1}{24} \left( \frac{1}{r^6} - \frac{12 \dot{r}^2}{r^5} + \frac{3 \ddot{r}}{r^4} \right) (δt)^4 + ... \label{13.8.13} \tag{13.8.13}\]
and \[G = δt - \frac{1}{6r^3} (δt)^3 + \frac{\dot{r}}{4r^4} (δt)^4 + ... \label{13.8.14} \tag{13.8.14}\]
Now we are going to look at the triangle and sector areas. From figure \(\text{XIII.1}\) we can see that
\[\textbf{A}_1 = \frac{1}{2} \textbf{r}_2 \times \textbf{r}_3 , \ \textbf{A}_2 = \frac{1}{2} \textbf{r}_1 \times \textbf{r}_3 , \ \textbf{A}_3 = \frac{1}{2} \textbf{r}_1 \times \textbf{r}_2 . \label{13.8.15a,b,c} \tag{13.8.15a,b,c}\]
Also, angular momentum per unit mass is \(\textbf{r} \times \textbf{v}\) and Kepler’s second law tells us that areal speed is half the angular momentum per unit mass and that it is constant and equal to \(\frac{1}{2} \sqrt{l}\) (in the units that we are using), so that
\[\dot{\textbf{B}}_1 = \frac{1}{2} \textbf{r}_1 \times \dot{\textbf{r}}_1 = \frac{1}{2} \textbf{r}_2 \times \dot{\textbf{r}}_2 = \frac{1}{2} \textbf{r}_3 \times \dot{\textbf{r}}_3 .\label{13.8.16a,b,c} \tag{13.8.16a,b,c}\]
All four of these vectors are parallel and perpendicular to the plane of the orbit, to that their magnitudes are just equal to their \(z\)-components. From the usual formulas for the components of a vector product we have, then,
\[A_1 = \frac{1}{2}(x_2 y_3 - y_2 x_3 ) , \quad A_2 = \frac{1}{2} (x_1 y_3 - y_1 x_3 ) , \quad A_3 = \frac{1}{2} (x_1 y_2 - y_1 x_2 ) \label{13.8.17a,b,c} \tag{13.8.17a,b,c}\]
and
\[\frac{1}{2} \sqrt{l} = \frac{1}{2} (x_1 \dot{y}_1 - y_1 \dot{x}_1 ) = \frac{1}{2}(x_2 \dot{y}_2 - y_2 \dot{x}_2 ) = \frac{1}{2} (x_3 \dot{y}_3 - y_3 \dot{x}_3 ) . \label{13.8.18a,b,c} \tag{13.8.18a,b,c}\]
Now, start from the second observation \((x_2 , y_2)\) at instant \(t_2\). We shall try to predict the third observation, using equations 13.8.11-14, in which \(x + δx\) is \(x_3\) and \(δt\) is \(t_3 − t_2\), which we are calling (see section 13.3) \(τ_1\). I shall make the subscripts for \(F\) and \(G\) the same as the subscripts for \(τ\). Thus the \(F\) and \(G\) that connect observations 2 and 3 will have subscript 1, just as we are using the notation \(τ_1\) for \(t_3 − t_2\).
Thus we have
\[x_2 = F_1 x_2 + G_1 \dot{x}_2 \label{13.8.19} \tag{13.8.19}\]
and \[y_3 = F_1 y_2 + G_1 \dot{y}_2 , \label{13.8.20} \tag{13.8.20}\]
where \[F_1 = 1 - \frac{1}{2r_2^3}τ_1^2 + \frac{\dot{r}_2}{2r_2^4} τ_1^3 + \frac{1}{24} \left( \frac{1}{r_2^6} - \frac{12 \dot{r}_2^2}{r_2^5} + \frac{3 \ddot{r}_2}{r_2^4} \right) τ_1^4 + ... \label{13.8.21} \tag{13.8.21}\]
and \[G_1 = τ_1 - \frac{1}{6r_2^3} τ_1^3 + \frac{\dot{r}_2}{4r_2^3} τ_1^4 + ... \quad . \label{13.8.22} \tag{13.8.22}\]
Similarly, the first observation is given by
\[x_1 = F_3 x_2 + G_3 \dot{x}_2 \label{13.8.23} \tag{13.8.23}\]
and \[y_1 = F_3 y_2 + G_3 \dot{y}_2 , \label{13.8.24} \tag{13.8.24}\]
where, by substitution of \(−τ_3\) for \(δt\),
\[F_3 = 1 - \frac{1}{2r_2^3} τ_3^2 - \frac{\dot{r}_2}{2r_2^4} τ_3^3 + \frac{1}{24} \left( \frac{1}{r_2^6} - \frac{12 \dot{r}_2^2}{r_2^5} + \frac{3\ddot{r}_2}{r_2^4} \right) τ_3^4 + ... \label{13.8.25} \tag{13.8.25}\]
and \[G_3 = -τ_3 + \frac{1}{6r_2^3} τ_3^3 + \frac{\dot{r}_2}{4r_2^4}τ_3^4 + ... \quad . \label{13.8.26} \tag{13.8.26}\]
From equations 13.8.17,18,19,20,23,24, we soon find that
\[A_1 = \frac{1}{2} G_1 \sqrt{l} , \quad A_2 = \frac{1}{2} ( F_3 G_1 - F_1G_3 ) \sqrt{l} , \quad A_3 = -\frac{1}{2} G_3 \sqrt{l} . \label{13.8.27a,b,c} \tag{13.8.27a,b,c}\]
Now we do not yet know \(\dot{r}\) or \(\ddot{r}\), but we can take the expansions of \(F\) and \(G\) as far as \(τ^2\) . We then obtain, correct to \(τ^3\) :
\[A_1 = \frac{1}{2} \sqrt{l} τ_1 \left( 1 - \frac{τ_1^2}{6_2^3} \right) , \label{13.8.28} \tag{13.8.28}\]
\[A_2 = \frac{1}{2} \sqrt{l} τ_2 \left( 1 - \frac{τ_1^2}{6r_2^3} \right) , \label{13.8.29} \tag{13,8.29}\]
and \[A_3 = \frac{1}{2} \sqrt{l} τ_3 \left( 1 - \frac{τ_3^2}{6r_2^3} \right) . \label{13.8.30} \tag{13.8.30}\]
Thus the triangle ratio \(a_1 = A_1/A_2\) is
\[a_1 = \frac{τ_1}{τ_2} \left( 1 - \frac{τ_1^2}{6r_2^3} \right) \left( 1 - \frac{τ_2^2}{6r_2^3} \right)^{-1} , \label{13.8.31} \tag{13.8.31}\]
which, to order \(τ^3\) , is \[a_1 = \frac{τ_1}{τ_2} \left( 1 + \frac{(τ_2^2 - τ_1^2)}{6r_2^3} \right) , \label{13.8.32} \tag{13.8.32}\]
or, with \(τ_2 - τ_1 = τ_3\), \[a_1 = \frac{τ_1}{τ_2} \left( 1 + \frac{τ_3 (τ_2 + τ_1)}{6r_2^3} \right) . \label{13.8.33} \tag{13.8.33}\]
Similarly, \[a_3 = \frac{τ_3}{τ_2} \left( 1 + \frac{τ_1 (τ_2 + τ_3)}{6r_2^3} \right) . \label{13.8.34} \tag{13.8.34}\]
Further, with \(τ_1 /τ_2 = b_1\) and \(τ_3 /τ_2 = b_3\),
\[a_1 = b_1 + \frac{τ_1 τ_3}{6r_2^3} (1 + b_1) \quad \text{and} \quad a_3 = b_3 + \frac{τ_1 τ_3}{6r_2^3} (1 + b_3) . \label{13.8.35a,b} \tag{13.8.35a,b}\]
These will serve as better approximations for the triangle ratios. Be aware, however, that equations 13.8.35a,b are only approximations, and do not give the exact values for \(a_1\) and \(a_3\).
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How to: Add and Remove Items from a ConcurrentDictionary
This example shows how to add, retrieve, update, and remove items from a System.Collections.Concurrent.ConcurrentDictionary<TKey,TValue>. This collection class is a thread-safe implementation. We recommend that you use it whenever multiple threads might be attempting to access the elements concurrently.
ConcurrentDictionary<TKey,TValue> provides several convenience methods that make it unnecessary for code to first check whether a key exists before it attempts to add or remove data. The following table lists these convenience methods and describes when to use them.
Method Use when…
AddOrUpdate You want to add a new value for a specified key and, if the key already exists, you want to replace its value.
GetOrAdd You want to retrieve the existing value for a specified key and, if the key does not exist, you want to specify a key/value pair.
TryAdd, TryGetValue , TryUpdate , TryRemove You want to add, get, update, or remove a key/value pair, and, if the key already exists or the attempt fails for any other reason, you want to take some alternative action.
Example
The following example uses two Task instances to add some elements to a ConcurrentDictionary<TKey,TValue> concurrently, and then outputs all of the contents to show that the elements were added successfully. The example also shows how to use the AddOrUpdate, TryGetValue, and GetOrAdd methods to add, update, and retrieve items from the collection.
namespace DictionaryHowTo
{
using System;
using System.Collections.Concurrent;
using System.Collections.Generic;
using System.Linq;
using System.Text;
using System.Threading;
using System.Threading.Tasks;
// The type of the Value to store in the dictionary:
class CityInfo : IEqualityComparer<CityInfo>
{
public string Name { get; set; }
public DateTime lastQueryDate { get; set; }
public decimal Longitude { get; set; }
public decimal Latitude { get; set; }
public int[] RecentHighTemperatures { get; set; }
public CityInfo(string name, decimal longitude, decimal latitude, int[] temps)
{
Name = name;
lastQueryDate = DateTime.Now;
Longitude = longitude;
Latitude = latitude;
RecentHighTemperatures = temps;
}
public CityInfo()
{
}
public CityInfo(string key)
{
Name = key;
// MaxValue means "not initialized"
Longitude = Decimal.MaxValue;
Latitude = Decimal.MaxValue;
lastQueryDate = DateTime.Now;
RecentHighTemperatures = new int[] { 0 };
}
public bool Equals(CityInfo x, CityInfo y)
{
return x.Name == y.Name && x.Longitude == y.Longitude && x.Latitude == y.Latitude;
}
public int GetHashCode(CityInfo obj)
{
CityInfo ci = (CityInfo)obj;
return ci.Name.GetHashCode();
}
}
class Program
{
// Create a new concurrent dictionary.
static ConcurrentDictionary<string, CityInfo> cities = new ConcurrentDictionary<string, CityInfo>();
static void Main(string[] args)
{
CityInfo[] data =
{
new CityInfo(){ Name = "Boston", Latitude = 42.358769M, Longitude = -71.057806M, RecentHighTemperatures = new int[] {56, 51, 52, 58, 65, 56,53}},
new CityInfo(){ Name = "Miami", Latitude = 25.780833M, Longitude = -80.195556M, RecentHighTemperatures = new int[] {86,87,88,87,85,85,86}},
new CityInfo(){ Name = "Los Angeles", Latitude = 34.05M, Longitude = -118.25M, RecentHighTemperatures = new int[] {67,68,69,73,79,78,78}},
new CityInfo(){ Name = "Seattle", Latitude = 47.609722M, Longitude = -122.333056M, RecentHighTemperatures = new int[] {49,50,53,47,52,52,51}},
new CityInfo(){ Name = "Toronto", Latitude = 43.716589M, Longitude = -79.340686M, RecentHighTemperatures = new int[] {53,57, 51,52,56,55,50}},
new CityInfo(){ Name = "Mexico City", Latitude = 19.432736M, Longitude = -99.133253M, RecentHighTemperatures = new int[] {72,68,73,77,76,74,73}},
new CityInfo(){ Name = "Rio de Janiero", Latitude = -22.908333M, Longitude = -43.196389M, RecentHighTemperatures = new int[] {72,68,73,82,84,78,84}},
new CityInfo(){ Name = "Quito", Latitude = -0.25M, Longitude = -78.583333M, RecentHighTemperatures = new int[] {71,69,70,66,65,64,61}}
};
// Add some key/value pairs from multiple threads.
Task[] tasks = new Task[2];
tasks[0] = Task.Run(() =>
{
for (int i = 0; i < 2; i++)
{
if (cities.TryAdd(data[i].Name, data[i]))
Console.WriteLine("Added {0} on thread {1}", data[i],
Thread.CurrentThread.ManagedThreadId);
else
Console.WriteLine("Could not add {0}", data[i]);
}
});
tasks[1] = Task.Run(() =>
{
for (int i = 2; i < data.Length; i++)
{
if (cities.TryAdd(data[i].Name, data[i]))
Console.WriteLine("Added {0} on thread {1}", data[i],
Thread.CurrentThread.ManagedThreadId);
else
Console.WriteLine("Could not add {0}", data[i]);
}
});
// Output results so far.
Task.WaitAll(tasks);
// Enumerate collection from the app main thread.
// Note that ConcurrentDictionary is the one concurrent collection
// that does not support thread-safe enumeration.
foreach (var city in cities)
{
Console.WriteLine("{0} has been added.", city.Key);
}
AddOrUpdateWithoutRetrieving();
RetrieveValueOrAdd();
RetrieveAndUpdateOrAdd();
Console.WriteLine("Press any key.");
Console.ReadKey();
}
// This method shows how to add key-value pairs to the dictionary
// in scenarios where the key might already exist.
private static void AddOrUpdateWithoutRetrieving()
{
// Sometime later. We receive new data from some source.
CityInfo ci = new CityInfo() { Name = "Toronto",
Latitude = 43.716589M,
Longitude = -79.340686M,
RecentHighTemperatures = new int[] { 54, 59, 67, 82, 87, 55, -14 } };
// Try to add data. If it doesn't exist, the object ci is added. If it does
// already exist, update existingVal according to the custom logic in the
// delegate.
cities.AddOrUpdate(ci.Name, ci,
(key, existingVal) =>
{
// If this delegate is invoked, then the key already exists.
// Here we make sure the city really is the same city we already have.
// (Support for multiple cities of the same name is left as an exercise for the reader.)
if (ci != existingVal)
throw new ArgumentException("Duplicate city names are not allowed: {0}.", ci.Name);
// The only updatable fields are the temperature array and lastQueryDate.
existingVal.lastQueryDate = DateTime.Now;
existingVal.RecentHighTemperatures = ci.RecentHighTemperatures;
return existingVal;
});
// Verify that the dictionary contains the new or updated data.
Console.Write("Most recent high temperatures for {0} are: ", cities[ci.Name].Name);
int[] temps = cities[ci.Name].RecentHighTemperatures;
foreach (var temp in temps) Console.Write("{0}, ", temp);
Console.WriteLine();
}
// This method shows how to use data and ensure that it has been
// added to the dictionary.
private static void RetrieveValueOrAdd()
{
string searchKey = "Caracas";
CityInfo retrievedValue = null;
try
{
retrievedValue = cities.GetOrAdd(searchKey, GetDataForCity(searchKey));
}
catch (ArgumentException e)
{
Console.WriteLine(e.Message);
}
// Use the data.
if (retrievedValue != null)
{
Console.Write("Most recent high temperatures for {0} are: ", retrievedValue.Name);
int[] temps = cities[retrievedValue.Name].RecentHighTemperatures;
foreach (var temp in temps) Console.Write("{0}, ", temp);
}
Console.WriteLine();
}
// This method shows how to retrieve a value from the dictionary,
// when you expect that the key/value pair already exists,
// and then possibly update the dictionary with a new value for the key.
private static void RetrieveAndUpdateOrAdd()
{
CityInfo retrievedValue;
string searchKey = "Buenos Aires";
if (cities.TryGetValue(searchKey, out retrievedValue))
{
// use the data
Console.Write("Most recent high temperatures for {0} are: ", retrievedValue.Name);
int[] temps = retrievedValue.RecentHighTemperatures;
foreach (var temp in temps) Console.Write("{0}, ", temp);
// Make a copy of the data. Our object will update its lastQueryDate automatically.
CityInfo newValue = new CityInfo(retrievedValue.Name,
retrievedValue.Longitude,
retrievedValue.Latitude,
retrievedValue.RecentHighTemperatures);
// Replace the old value with the new value.
if (!cities.TryUpdate(searchKey, newValue, retrievedValue))
{
//The data was not updated. Log error, throw exception, etc.
Console.WriteLine("Could not update {0}", retrievedValue.Name);
}
}
else
{
// Add the new key and value. Here we call a method to retrieve
// the data. Another option is to add a default value here and
// update with real data later on some other thread.
CityInfo newValue = GetDataForCity(searchKey);
if(cities.TryAdd(searchKey, newValue))
{
// use the data
Console.Write("Most recent high temperatures for {0} are: ", newValue.Name);
int[] temps = newValue.RecentHighTemperatures;
foreach (var temp in temps) Console.Write("{0}, ", temp);
}
else
Console.WriteLine("Unable to add data for {0}", searchKey);
}
}
//Assume this method knows how to find long/lat/temp info for any specified city.
static CityInfo GetDataForCity(string name)
{
// Real implementation left as exercise for the reader.
if (String.CompareOrdinal(name, "Caracas") == 0)
return new CityInfo() { Name = "Caracas",
Longitude = 10.5M,
Latitude = -66.916667M,
RecentHighTemperatures = new int[] { 91, 89, 91, 91, 87, 90, 91 } };
else if (String.CompareOrdinal(name, "Buenos Aires") == 0)
return new CityInfo() { Name = "Buenos Aires",
Longitude = -34.61M,
Latitude = -58.369997M,
RecentHighTemperatures = new int[] { 80, 86, 89, 91, 84, 86, 88 } };
else
throw new ArgumentException("Cannot find any data for {0}", name);
}
}
}
Imports System
Imports System.Collections.Concurrent
Imports System.Collections.Generic
Imports System.Linq
Imports System.Text
Imports System.Threading
Imports System.Threading.Tasks
Namespace DictionaryHowToVB
' The type of the value to store in the dictionary
Class CityInfo
Private _name As String
Property Name As String
Get
Return _name
End Get
Set(ByVal value As String)
_name = value
End Set
End Property
Private _lastQueryDate As DateTime
Property LastQueryDate As DateTime
Get
Return _lastQueryDate
End Get
Set(ByVal value As DateTime)
_lastQueryDate = value
End Set
End Property
Private _longitude As Decimal
Property Longitude As Decimal
Get
Return _longitude
End Get
Set(ByVal value As Decimal)
_longitude = value
End Set
End Property
Private _latitude As Decimal
Property Latitude As Decimal
Get
Return _latitude
End Get
Set(ByVal value As Decimal)
_latitude = value
End Set
End Property
Private _highTemps() As Integer
Property RecentHighTemperatures As Integer()
Get
Return _highTemps
End Get
Set(ByVal value As Integer())
_highTemps = value
End Set
End Property
Public Sub New()
End Sub
Public Sub New(ByVal key As String)
_name = key
' MaxValue means "not initialized"
_longitude = Decimal.MaxValue
_latitude = Decimal.MaxValue
_lastQueryDate = DateTime.Now
_highTemps = {0}
End Sub
Public Sub New(ByVal name As String, ByVal longitude As Decimal,
ByVal latitude As Decimal, ByVal temps As Integer())
_name = name
_longitude = longitude
_latitude = latitude
_highTemps = temps
End Sub
End Class
Class Program
' Create a new concurrent dictionary with the specified concurrency level and capacity.
Shared cities As New ConcurrentDictionary(Of String, CityInfo)(System.Environment.ProcessorCount, 10)
Shared Sub Main()
Dim data As CityInfo() =
{New CityInfo With {.Name = "Boston", .Latitude = 42.358769, .Longitude = -71.057806, .RecentHighTemperatures = {56, 51, 52, 58, 65, 56, 53}},
New CityInfo With {.Name = "Miami", .Latitude = 25.780833, .Longitude = -80.195556, .RecentHighTemperatures = {86, 87, 88, 87, 85, 85, 86}},
New CityInfo With {.Name = "Los Angeles", .Latitude = 34.05, .Longitude = -118.25, .RecentHighTemperatures = {67, 68, 69, 73, 79, 78, 78}},
New CityInfo With {.Name = "Seattle", .Latitude = 47.609722, .Longitude = -122.333056, .RecentHighTemperatures = {49, 50, 53, 47, 52, 52, 51}},
New CityInfo With {.Name = "Toronto", .Latitude = 43.716589, .Longitude = -79.340686, .RecentHighTemperatures = {53, 57, 51, 52, 56, 55, 50}},
New CityInfo With {.Name = "Mexico City", .Latitude = 19.432736, .Longitude = -99.133253, .RecentHighTemperatures = {72, 68, 73, 77, 76, 74, 73}},
New CityInfo With {.Name = "Rio de Janiero", .Latitude = -22.908333, .Longitude = -43.196389, .RecentHighTemperatures = {72, 68, 73, 82, 84, 78, 84}},
New CityInfo With {.Name = "Quito", .Latitude = -0.25, .Longitude = -78.583333, .RecentHighTemperatures = {71, 69, 70, 66, 65, 64, 61}}}
' Add some key/value pairs from multiple threads.
Dim tasks(1) As Task
tasks(0) = Task.Run(Sub()
For i As Integer = 0 To 1
If cities.TryAdd(data(i).Name, data(i)) Then
Console.WriteLine("Added {0} on thread {1}", data(i).Name, Thread.CurrentThread.ManagedThreadId)
Else
Console.WriteLine("Could not add {0}", data(i))
End If
Next
End Sub)
tasks(1) = Task.Run(Sub()
For i As Integer = 2 To data.Length - 1
If cities.TryAdd(data(i).Name, data(i)) Then
Console.WriteLine("Added {0} on thread {1}", data(i).Name, Thread.CurrentThread.ManagedThreadId)
Else
Console.WriteLine("Could not add {0}", data(i))
End If
Next
End Sub)
' Output results so far
Task.WaitAll(tasks)
' Enumerate data on main thread. Note that
' ConcurrentDictionary is the one collection class
' that does not support thread-safe enumeration.
For Each city In cities
Console.WriteLine("{0} has been added", city.Key)
Next
AddOrUpdateWithoutRetrieving()
RetrieveValueOrAdd()
RetrieveAndUpdateOrAdd()
Console.WriteLine("Press any key")
Console.ReadKey()
End Sub
' This method shows how to add key-value pairs to the dictionary
' in scenarios where the key might already exist.
Private Shared Sub AddOrUpdateWithoutRetrieving()
' Sometime later. We receive new data from some source.
Dim ci = New CityInfo With {.Name = "Toronto", .Latitude = 43.716589, .Longitude = -79.340686, .RecentHighTemperatures = {54, 59, 67, 82, 87, 55, -14}}
' Try to add data. If it doesn't exist, the object ci is added. If it does
' already exist, update existingVal according to the custom logic in the
' delegate.
cities.AddOrUpdate(ci.Name, ci, Function(key, existingVal)
' If this delegate is invoked, then the key already exists.
' Here we make sure the city really is the same city we already have.
' (Support for multiple keys of the same name is left as an exercise for the reader.)
If (ci.Name = existingVal.Name And ci.Longitude = existingVal.Longitude) = False Then
Throw New ArgumentException("Duplicate city names are not allowed: {0}.", ci.Name)
End If
' The only updatable fields are the temperature array and lastQueryDate.
existingVal.LastQueryDate = DateTime.Now
existingVal.RecentHighTemperatures = ci.RecentHighTemperatures
Return existingVal
End Function)
' Verify that the dictionary contains the new or updated data.
Console.Write("Most recent high temperatures for {0} are: ", cities(ci.Name).Name)
Dim temps = cities(ci.Name).RecentHighTemperatures
For Each temp In temps
Console.Write("{0}, ", temp)
Next
Console.WriteLine()
End Sub
'This method shows how to use data and ensure that it has been
' added to the dictionary.
Private Shared Sub RetrieveValueOrAdd()
Dim searchKey = "Caracas"
Dim retrievedValue As CityInfo = Nothing
Try
retrievedValue = cities.GetOrAdd(searchKey, GetDataForCity(searchKey))
Catch e As ArgumentException
Console.WriteLine(e.Message)
End Try
' Use the data.
If Not retrievedValue Is Nothing Then
Console.WriteLine("Most recent high temperatures for {0} are: ", retrievedValue.Name)
Dim temps = cities(retrievedValue.Name).RecentHighTemperatures
For Each temp In temps
Console.Write("{0}, ", temp)
Next
End If
Console.WriteLine()
End Sub
' This method shows how to retrieve a value from the dictionary,
' when you expect that the key/value pair already exists,
' and then possibly update the dictionary with a new value for the key.
Private Shared Sub RetrieveAndUpdateOrAdd()
Dim retrievedValue As CityInfo = New CityInfo()
Dim searchKey = "Buenos Aires"
If (cities.TryGetValue(searchKey, retrievedValue)) Then
' Use the data
Console.Write("Most recent high temperatures for {0} are: ", retrievedValue.Name)
Dim temps = retrievedValue.RecentHighTemperatures
For Each temp In temps
Console.Write("{0}, ", temp)
Next
' Make a copy of the data. Our object will update its lastQueryDate automatically.
Dim newValue As CityInfo = New CityInfo(retrievedValue.Name,
retrievedValue.Longitude,
retrievedValue.Latitude,
retrievedValue.RecentHighTemperatures)
Else
Console.WriteLine("Unable to find data for {0}", searchKey)
End If
End Sub
' Assume this method knows how to find long/lat/temp info for any specified city.
Private Shared Function GetDataForCity(ByVal searchKey As String) As CityInfo
' Real implementation left as exercise for the reader.
If String.CompareOrdinal(searchKey, "Caracas") = 0 Then
Return New CityInfo() With {.Name = "Caracas",
.Longitude = 10.5,
.Latitude = -66.916667,
.RecentHighTemperatures = {91, 89, 91, 91, 87, 90, 91}}
ElseIf String.CompareOrdinal(searchKey, "Buenos Aires") = 0 Then
Return New CityInfo() With {.Name = "Buenos Aires",
.Longitude = -34.61,
.Latitude = -58.369997,
.RecentHighTemperatures = {80, 86, 89, 91, 84, 86, 88}}
Else
Throw New ArgumentException("Cannot find any data for {0}", searchKey)
End If
End Function
End Class
End Namespace
ConcurrentDictionary<TKey,TValue> is designed for multithreaded scenarios. You do not have to use locks in your code to add or remove items from the collection. However, it is always possible for one thread to retrieve a value, and another thread to immediately update the collection by giving the same key a new value.
Also, although all methods of ConcurrentDictionary<TKey,TValue> are thread-safe, not all methods are atomic, specifically GetOrAdd and AddOrUpdate. The user delegate that is passed to these methods is invoked outside of the dictionary's internal lock (this is done to prevent unknown code from blocking all threads). Therefore, it is possible for this sequence of events to occur:
1) threadA calls GetOrAdd, finds no item and creates a new item to Add by invoking the valueFactory delegate.
2) threadB calls GetOrAdd concurrently, its valueFactory delegate is invoked and it arrives at the internal lock before threadA, and so its new key-value pair is added to the dictionary.
3) threadA's user delegate completes, and the thread arrives at the lock, but now sees that the item exists already.
4) threadA performs a "Get" and returns the data that was previously added by threadB.
Therefore, it is not guaranteed that the data that is returned by GetOrAdd is the same data that was created by the thread's valueFactory. A similar sequence of events can occur when AddOrUpdate is called.
See also
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| 0.9718 |
Towards a risk evaluation of workers' exposure to handborne and airborne microbial species as exemplified with waste collection workers
Tidsskriftartikel - 2020
Resume
Bioaerosol exposure is associated with health problems. The aim of this study is to evaluate whether it is possible to assess the risks posed by waste collection workers' exposure through identification and characterization of bacterial and fungal species, to which the workers are exposed. Using MALDI-TOF MS, microorganisms in waste collection workers' exposure through air, hand, and contact with the steering wheel were identified. Fungi found in high concentrations from the workers' exposure were characterized for the total inflammatory potential (TIP), cytotoxicity, and biofilm-forming capacity. In total, 180 different bacterial and 37 different fungal species in the workers' exposure samples were identified. Some of them belong to Risk Group 2, e.g. Escherichia coli, Klebsiella oxytoca, Staphylococcus aureus, and Aspergillus fumigatus, some have been associated with occupational health problems e.g. Penicillium citrinum and P. glabrum and some are described as emerging pathogens e.g. Aureobasidium pullulans. The TIP of fungal species was dose-dependent. High TIP values were found for Penicillium italicum, P. brevicompactum, P. citrinum, and P. glabrum. Several species were cytotoxic, e.g. A. niger and P. expansum, while some, e.g. P. chrysogenum, did not affect the cell viability. Based on waste workers' average inhalation rate, they inhaled up to 2.3 × 104 cfu of A. niger, 7.4 × 104 cfu of P. expansum, and 4.0 × 106 cfu of P. italicum per work day. Some species e.g. A. niger and P. citrinum were able to form biofilm. In conclusion, the workers were exposed to several species of microorganisms of which some to varying degrees can be evaluated concerning risk. Thus, some microorganisms belong to Risk Group 2, and some are described as causing agents of occupational health problems, emerging pathogens, or intrinsically antibiotic resistant. For some other species very little is known. The TIP, cytotoxicity, and ability to form biofilm of the dominating fungi support and expand previous findings. These parameters depended on the species and the dose, thus highlighting the importance of species identification and exposure level in the risk assessment of exposure.
Reference
Madsen AM, Frederiksen MW, Jacobsen MH, Tendal K. Towards a risk evaluation of workers' exposure to handborne and airborne microbial species as exemplified with waste collection workers. Environmental Research 2020;183:109177.
doi: 10.1016/j.envres.2020.109177
Gå til Tidsskriftartikel
|
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| 0.803117 |
@Article{acp-21-3059-2021, AUTHOR = {Wang, G. and Luo, Z. and Huang, Y. and Sun, W. and Wei, Y. and Xiao, L. and Deng, X. and Zhu, J. and Li, T. and Zhang, W.}, TITLE = {Simulating the spatiotemporal variations in aboveground biomass in Inner Mongolian grasslands under environmental changes}, JOURNAL = {Atmospheric Chemistry and Physics}, VOLUME = {21}, YEAR = {2021}, NUMBER = {4}, PAGES = {3059--3071}, URL = {https://acp.copernicus.org/articles/21/3059/2021/}, DOI = {10.5194/acp-21-3059-2021} }
|
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| 0.995244 |
Popularity
6.6
Stable
Activity
3.9
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1,603
115
489
Description
## Features of PSpider - Support multi-threading crawling mode (using threading and requests) - Support distributed crawling mode (using threading, requests and redis) - Support using proxies for crawling (using threading and queue) - Define some utility functions and classes, for example: UrlFilter, get_string_num, etc - Fewer lines of code, easyer to read, understand and expand
## Modules of PSpider - utilities module: define some utilities functions and classes for spider - instances module: define classes of fetcher, parser, saver for multi-threading spider - concurrent module: define WebSpiderFrame of multi-threading spider and distributed spider
Programming language: Python
License: BSD 2-clause "Simplified" License
PSpider alternatives and similar packages
Based on the "Web Crawling" category.
Alternatively, view PSpider alternatives based on common mentions on social networks and blogs.
Do you think we are missing an alternative of PSpider or a related project?
Add another 'Web Crawling' Package
README
PSpider
A simple web spider frame written by Python, which needs Python3.5+
Features of PSpider
1. Support multi-threading crawling mode (using threading and requests)
2. Support multi-processing in parse process, automatically (using multiprocessing)
3. Support using proxies for crawling (using threading and queue)
4. Define some utility functions and classes, for example: UrlFilter, get_string_num, etc
5. Fewer lines of code, easyer to read, understand and expand
Modules of PSpider
1. utilities module: define some utilities functions and classes for multi-threading spider
2. instances module: define classes of Fetcher, Parser, Saver for multi-threading spider
3. concurrent module: define WebSpiderFrame of multi-threading spider
Procedure of PSpider
[](procedure.png) ①: Fetchers get url from UrlQueue, and make requests based on this url
②: Put the result of ① to HtmlQueue, and so Parser can get it
③: Parser gets item from HtmlQueue, and parses it to get new urls and items
④: Put the new urls to UrlQueue, and so Fetcher can get it
⑤: Put the items to ItemQueue, and so Saver can get it
⑥: Saver gets item from ItemQueue, and saves it to filesystem or database
⑦: Proxieser gets proxies from web or database, and puts proxies to ProxiesQueue
⑧: Fetcher gets proxies from ProxiesQueue if needed, and makes requests based on this proxies
Tutorials of PSpider
Installation: you'd better use the first method
（1）Copy the "spider" directory to your project directory, and import spider
（2）Install spider to your python system using python3 setup.py install
See test.py
TodoList
1. More Demos
2. Distribute Spider
3. Execute JavaScript
If you have any questions or advices, you can commit "Issues" or "Pull requests"
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| 0.991854 |
Translate
Image of Advanced Programming in the UNIX Environment, Second Edition (Addison-Wesley Professional Computing Series)
Image of Beginning Google Maps API 3
Image of Operating System Concepts
Image of Modern Operating Systems (3rd Edition)
More GNOME Shell Customization
If you plan to customize the GNOME Shell in any meaningful way, you need to understand the technologies underlying the GNOME Shell and understand how to write a GNOME Shell extension to provide the customization that you require. In this post I delve deeper into the technologies behind the new GNOME Shell and provide sample code for a number of simple extensions which demonstrate how to customize and extend various components of the GNOME Shell user interface.
Essentially, the GNOME Shell is an integrated window manager, compositor, and application launcher. It acts as a composting window manager for the desktop displaying both application windows and other objects in a Clutter scene graph. Most of the UI code is written in JavaScript which accesses Clutter and other underlying libraries via GObject Introspection.
Here is a block diagram of the underlying technologies that support the GNOME Shell as of v3.0:
GNOME3 Shell screenshot
This is a modified version of a diagram that exists on the GNOME website.
The GNOME Shell uses OpenGL to render graphics. OpenGL uses a hardware accelerated pixel format by default but can support software rendering. However, hardware acceleration is required to run the GNOME Shell as it uses a number of 3D capabilities to accelerate the transforms. Most graphics cards less than 3 years old should support hardware acceleration. If hardware acceleration is unavailable, the GNOME Shell defaults back to a modified version of the GNOME 2 Panel. See Vincent Untz’s post for further information on this fallback mode. In addition, you can force the GNOME Shell to use the fallback mode via a switch in the Settings, System Info panel. Whether a graphics card supports hardware acceleration or not is determined by a helper application /usr/libexec/gnome-session-is-accelerated which returns either 0 or 1. Another limitation in the 3.0 version of the GNOME shell is that the total horizontal or vertical pixel count of all your monitors must be less than 2048.
int
main (int argc, char **argv)
{
Display *display = NULL;
int ret = 1;
display = XOpenDisplay (NULL);
if (!display) {
_print_error ("No X display.");
goto out;
}
if (_has_composite (display) != 0) {
_print_error ("No composite extension.");
goto out;
}
if (_has_hardware_gl (display) != 0) {
_print_error ("No hardware 3D support.");
goto out;
}
if (_has_texture_from_pixmap (display) != 0) {
_print_error ("No GLX_EXT_texture_from_pixmap support.");
goto out;
}
if (_is_max_texture_size_big_enough (display) != 0) {
_print_error ("GL_MAX_TEXTURE_SIZE is too small.");
goto out;
}
ret = 0;
out:
if (display)
XCloseDisplay (display);
return ret;
}
As you can see, the use of gotos is alive and well! If you are certain that your graphics card supports hardware acceleration but for some reason GNOME Shell only works in fallback mode, try replacing this application with a script that simply returns true. There is also a black list of graphic cards. According to Emmanuele Bassi:
the hard limit is on Intel 915/945 integrated graghics (which are sadly what netbooks generally use). The limit also used to exist on 965 but it was a soft limit that was lifted by fixing the driver in Mesa
You can also force the fallback mode using a gsettings key:
$ gsettings set org.gnome.desktop.session session-name gnome-fallback
Access to OpenGL is via Cogl which is a graphics API that exposes the features of 3D graphics hardware using a more object oriented design than OpenGL. The Clutter graphics library handles scene graphing. In Clutter, widgets are called actors, and windows are called stages. A Clutter application contains at least one stage containing actors such as rectangles, images, or text. A useful online resource for Clutter programming is Programming in Clutter by Murray Cumming. By the way, the Clutter library is also used in Moblin which, along with Maemo, is now part of Meego. Meego uses MX widgets on top of Clutter (a useful tutorial can be found here) whereas the GNOME Shell uses a Shell Toolkit (St) which implements many custom actors, such as containers, bins, boxes, buttons and scrollbars that are useful in implementing GNOME Shell UI features. The Shell Toolkit was derived from the Moblin UI Toolkit. See ../src/st in the gnome-shell GIT source code repository. The Shell Toolkit also implements CSS support (see ../src/st/st-theme*) which makes the GNOME Shell themeable. Generally if you see any object whole name starts with St. you can assume you are dealing with the Shell Toolkit. Accessibility is handled by Cally (Clutter Accessibility Implementation Library).
Window management is handled by a modified version of Metacity called Mutter. Before the introduction of Metacity in GNOME 2.2, GNOME used Enlightenment and then Sawfish as its window manager. Metacity uses the GTK+ graphical widget toolkit to create its user interface components, enabling it to be themeable and blend in with other GTK+ applications. Mutter is a newer compositing window manager based on Metacity and Clutter. Note that the GNOME Shell fallback mode still uses Metacity and Gtk+, whereas the normal hardware accelerated mode uses Mutter. By the way, the GNOME Shell is actually implemented as a Mutter plugin. You can obtain copious amounts of debug information from Mutter using imports.gi.Meta.set_verbose(true) or you can control the volume of debug information using the add_verbose_topic method, e.g. to view only window state use Meta = imports.gi.Meta; Meta.add_verbose_topic(Meta.DebugTopic.WINDOW_STATE).
The GObject Introspection layer sits on top of Mutter and the Shell toolkit. One way to look at this layer is to consider it a glue layer between the Mutter and Shell Tookit libraries and JavaScript. GObject Introspection is used to automatically generate the GObject bindings for JavaScript (gjs) which is the language used to implement the GNOME Shell UI. The actual version of the JavaScript language available using gjs is 1.8.5 as this JavaScript engine is based on Mozilla’s Spidermonkey JavaScript engine.
The goal of GObject Introspection is to describe the set of APIs in a uniform, machine readable XML format called GIR. A typelib is a compiled version of a GIR which is designed to be fast, memory efficient and complete enough so that language bindings can be written on top of it without other sources of information. You can examine the contents of a specific typelib file using g-ir-generate. For example, here is what is stored in the Shell Toolkit typelib file for st_texture_cache_load_uri_sync which I use in Example 7 below.
# g-ir-generate /usr/lib64/gnome-shell/St-1.0.typelib
....
<method name="load_uri_sync" c:identifier="st_texture_cache_load_uri_sync" throws="1">
<return-value transfer-ownership="none">
<type name="Clutter.Actor"/>
</return-value>
<parameters>
<parameter name="policy" transfer-ownership="none">
<type name="TextureCachePolicy"/>
</parameter>
<parameter name="uri" transfer-ownership="none">
<type name="utf8"/>
</parameter>
<parameter name="available_width" transfer-ownership="none">
<type name="gint32"/>
</parameter>
<parameter name="available_height" transfer-ownership="none">
<type name="gint32"/>
</parameter>
</parameters>
</method>
....
You also need to understand the various components that make up the GNOME Shell UI if you are going to be successful in customizing the GNOME Shell. Here are the various components of the screen displayed just after you successfully log in.
GNOME3 Shell screenshot
Here is a description of the various components:
• Top Bar (Also called Panel) – Horizontal bar at the top of the scrren. This is main access point to the shell. (../js/ui/panel.js)
• Activities button (Also called Hotspot) – Button and Hot Corner that brings up the Overview (see below). (../js/ui/panel.js)
• Application menu – Shows the name of the currently-focused application. You can click on it to get a menu of relevant actions. (../js/ui/panel.js)
• Clock – Also contains the Calendar (../js/ui/dateMenu.js)
• System Status Icons – Icons for accessibility options, Bluetooth, keyboard layout, etc. (../js/ui/status/*)
• Status Menu – Also contains the user menu. Lets you change your IM status, personal preferences, and exit the session. (../js/ui/statusMenu.js)
• Notifications Area – The part of the ephemeral horizontal bar at the bottom of screen where notifications are displayed. (../js/ui/notificationDaemon.js)
• Message tray – The part of the ephemeral horizontal bar at the bottom of screen where pending notifications and other messages display. (../js/ui/messageTray.js)
All the JavaScript code referenced above and in the next section is under /usr/share/gnome-shell.
The Overview screen is what you see when you click on the Activities button. It is mainly implemented in ../js/ui/overview.js. It has the following UI components:
GNOME3 Shell screenshot
Here is a description of the various components in this particular screen:
• Dash – Vertical bar on the left, that shows your favourite applications. (../js/ui/dash.js)
• View Selector – Lets you pick between Windows and Applications. (../js/ui/viewSelector.js, ../js/ui/overview.js)
• Search Entry – When you enter a string, various things (application names, documents, etc.) get searched. (../js/ui/viewSelector.js for the entry, and ../js/ui/searchDisplay.js for the display of search results)
• Workspace List – Vertical bar on the right, with thumbnails for the active workspaces. (../js/ui/workspaceSwitcherPopup.js)
In the following examples, I demonstrate how to customize various components of the GNOME Shell UI using extensions or by directly modifying the source code as in Example 7 if an extension is not possible. I assume that you know JavaScript and the components that form a GNOME Shell extension.
Example 1:
The GNOME Shell developers are pushing hard to eliminate the traditional notification area on the top bar of GNOME desktops. However for the moment, tray icons are still displayed on the Panel to the left of the System Status area.
For example gnote normally displays in the GNOME Shell message area as shown below:
GNOME3 Shell screenshot
With this simple extension:
const Panel = imports.ui.panel;
const StatusIconDispatcher = imports.ui.statusIconDispatcher;
function main() {
// add the notification(s) you want display on the top bar
// - one per line. Use the english text string displayed when
// hovering your mouse over the bottom right notification area
StatusIconDispatcher.STANDARD_TRAY_ICON_IMPLEMENTATIONS['gnote'] = 'gnote';
}
you can get gnote to display on the Panel.
GNOME3 Shell screenshot
Note how you import a module using the imports keyword. If you want to import a specific API version of a module, you can do so by specifying the required version number, e.g.
const Gtk = imports.gi.versions.Gtk = '3.0';
Example 2:
There is an existing extension in the gnome-shell-extensions source code repository called alternative status menu which replaces the GNOME Shell status menu in its entirety just so that the Power Off, Suspend and Hibernate menu options can be displayed instead of just the Suspend menu option.
GNOME3 Shell screenshot
Instead of the brute force approach used by this extension, the code shown below simply locates the suspend menu option and replaces it with the three required menu options and the auxiliary support functions.
const St = imports.gi.St;
const Main = imports.ui.main;
const PopupMenu = imports.ui.popupMenu;
const Shell = imports.gi.Shell;
const Lang = imports.lang;
const Gettext = imports.gettext.domain('gnome-shell');
const _ = Gettext.gettext;
function updateSuspend(object, pspec, item) {
item.actor.visible = object.get_can_suspend();
}
function updateHibernate(object, pspec, item) {
item.actor.visible = object.get_can_hibernate();
}
function onSuspendActivate(item) {
Main.overview.hide();
this._screenSaverProxy.LockRemote(Lang.bind(this, function() {
this._upClient.suspend_sync(null);
}));
}
function onHibernateActivate(item) {
Main.overview.hide();
this._screenSaverProxy.LockRemote(Lang.bind(this, function() {
this._upClient.hibernate_sync(null);
}));
}
function changeUserMenu()
{
let children = this.menu._getMenuItems();
for (let i = 0; i < children.length; i++) {
let item = children[i];
if (item.label) {
let _label = item.label.get_text();
// global.log("menu label: " + _label);
if (_label == _("Suspend"))
item.destroy();
}
}
let item = new PopupMenu.PopupMenuItem(_("Suspend"));
item.connect('activate', Lang.bind(this, onSuspendActivate));
this._upClient.connect('notify::can-suspend', Lang.bind(this, updateSuspend, item));
updateSuspend(this._upClient, null, item);
this.menu.addMenuItem(item);
item = new PopupMenu.PopupMenuItem(_("Hibernate"));
item.connect('activate', Lang.bind(this, onHibernateActivate));
this._upClient.connect('notify::can-hibernate', Lang.bind(this, updateHibernate, item));
updateHibernate(this._upClient, null, item);
this.menu.addMenuItem(item);
item = new PopupMenu.PopupMenuItem(_("Power Off..."));
item.connect('activate', Lang.bind(this, function() {
this._session.ShutdownRemote();
}));
this.menu.addMenuItem(item);
}
function main(metadata) {
// Post 3.0 let statusMenu = Main.panel._userMenu;
let statusMenu = Main.panel._statusmenu;
changeUserMenu.call(statusMenu);
}
There are different types of power hibernation but the above example just uses the default method. Some people might find it useful to have a sleep menu option also.
Note that I have commented out a line of code in the main function. The commented out line is what you should use in post 3.0 versions of the GNOME Shell. How you access Panel objects is changing. See GNOME Bugzilla 646915 for full details. Essentially, a number of Panel objects have been renamed and a _statusArea object now points to status area PanelButton objects. The idea is that you will be able to address each Panel object consistently as follows:
Main.panel._activities
Main.panel._appMenu
Main.panel._dateMenu
Main.panel._statusArea.a11y
Main.panel._statusArea.volume
Main.panel._userMenu
Example 3:
In this example. I show you how to add a menu to the middle of the Panel as shown below:
GNOME3 Shell screenshot
Here is the source code for the extension used to create that menu:
const St = imports.gi.St;
const Main = imports.ui.main;
const PopupMenu = imports.ui.popupMenu;
const PanelMenu = imports.ui.panelMenu;
const Gettext = imports.gettext;
const _ = Gettext.gettext;
function PlacesButton() {
this._init();
}
PlacesButton.prototype = {
__proto__: PanelMenu.Button.prototype,
_init: function() {
PanelMenu.Button.prototype._init.call(this, 0.0);
this._label = new St.Label({ text: _("MyPlaces") });
this.actor.set_child(this._label);
Main.panel._centerBox.add(this.actor, { y_fill: true });
let placeid;
this.placeItems = [];
this.defaultPlaces = Main.placesManager.getDefaultPlaces();
this.bookmarks = Main.placesManager.getBookmarks();
this.mounts = Main.placesManager.getMounts();
// Display default places
for ( placeid = 0; placeid < this.defaultPlaces.length; placeid++) {
this.placeItems[placeid] = new PopupMenu.PopupMenuItem(_(this.defaultPlaces[placeid].name));
this.placeItems[placeid].place = this.defaultPlaces[placeid];
this.menu.addMenuItem(this.placeItems[placeid]);
this.placeItems[placeid].connect('activate', function(actor,event) {
actor.place.launch();
});
}
this.menu.addMenuItem(new PopupMenu.PopupSeparatorMenuItem());
// Display default bookmarks
for ( let bookmarkid = 0; bookmarkid < this.bookmarks.length; bookmarkid++, placeid++) {
this.placeItems[placeid] = new PopupMenu.PopupMenuItem(_(this.bookmarks[bookmarkid].name));
this.placeItems[placeid].place = this.bookmarks[bookmarkid];
this.menu.addMenuItem(this.placeItems[placeid]);
this.placeItems[placeid].connect('activate', function(actor,event) {
actor.place.launch();
});
}
if (this.mounts.length > 0) {
this.menu.addMenuItem(new PopupMenu.PopupSeparatorMenuItem());
}
// Display default mounts
for ( let mountid = 0; mountid < this.mounts.length; placeid++, mountid++ ) {
this.placeItems[placeid] = new PopupMenu.PopupMenuItem(_(this.mounts[mountid].name));
this.placeItems[placeid].place = this.mounts[mountid];
this.menu.addMenuItem(this.placeItems[placeid]);
this.placeItems[placeid].connect('activate', function(actor,event) {
actor.place.launch();
});
}
Main.panel._centerBox.add(this.actor, { y_fill: true });
Main.panel._menus.addMenu(this.menu);
}
};
function main(extensionMeta) {
let userExtensionLocalePath = extensionMeta.path + '/locale';
Gettext.bindtextdomain("places_button", userExtensionLocalePath);
Gettext.textdomain("places_button");
new PlacesButton();
}
Notice how you can retrieve details of all places, bookmarks and mounts from Main.placesManager:
places = Main.placesManager.getDefaultPlaces();
bookmarks = Main.placesManager.getBookmarks();
mounts = Main.placesManager.getMounts();
Example 4:
In this example, I show you how to extend the previous example to display icons on each menu option as shown below:
GNOME3 Shell screenshot
Here is the modified source code:
const St = imports.gi.St;
const Main = imports.ui.main;
const PopupMenu = imports.ui.popupMenu;
const PanelMenu = imports.ui.panelMenu;
const Gettext = imports.gettext;
const _ = Gettext.gettext;
const MYPLACES_ICON_SIZE = 22;
function PlacesButton() {
this._init();
}
PlacesButton.prototype = {
__proto__: PanelMenu.Button.prototype,
_init: function() {
PanelMenu.Button.prototype._init.call(this, 0.0);
this._label = new St.Label({ text: _("MyPlaces") });
this.actor.set_child(this._label);
Main.panel._centerBox.add(this.actor, { y_fill: true });
let placeid;
this.placeItems = [];
this.defaultPlaces = Main.placesManager.getDefaultPlaces();
this.bookmarks = Main.placesManager.getBookmarks();
this.mounts = Main.placesManager.getMounts();
// Display default places
for ( placeid = 0; placeid < this.defaultPlaces.length; placeid++) {
this.placeItems[placeid] = new PopupMenu.PopupMenuItem(_(this.defaultPlaces[placeid].name));
let icon = this.defaultPlaces[placeid].iconFactory(MYPLACES_ICON_SIZE);
this.placeItems[placeid].addActor(icon, { align: St.Align.END});
this.menu.addMenuItem(this.placeItems[placeid]);
this.placeItems[placeid].connect('activate', function(actor,event) {
actor.place.launch();
});
}
this.menu.addMenuItem(new PopupMenu.PopupSeparatorMenuItem());
this.menu.addMenuItem(new PopupMenu.PopupSeparatorMenuItem());
// Display default bookmarks
for ( let bookmarkid = 0; bookmarkid < this.bookmarks.length; bookmarkid++, placeid++) {
this.placeItems[placeid] = new PopupMenu.PopupMenuItem(_(this.bookmarks[bookmarkid].name));
let icon = this.bookmarks[bookmarkid].iconFactory(MYPLACES_ICON_SIZE);
this.placeItems[placeid].addActor(icon, { align: St.Align.END});
this.menu.addMenuItem(this.placeItems[placeid]);
this.placeItems[placeid].connect('activate', function(actor,event) {
actor.place.launch();
});
}
if (this.mounts.length > 0) {
this.menu.addMenuItem(new PopupMenu.PopupSeparatorMenuItem());
}
// Display default mounts
for ( let mountid = 0; mountid < this.mounts.length; placeid++, mountid++ ) {
this.placeItems[placeid] = new PopupMenu.PopupMenuItem(_(this.mounts[mountid].name));
let icon = this.mounts[mountid].iconFactory(MYPLACES_ICON_SIZE);
this.placeItems[placeid].addActor(icon, { align: St.Align.END});
this.placeItems[placeid].place = this.mounts[mountid];
this.menu.addMenuItem(this.placeItems[placeid]);
this.placeItems[placeid].connect('activate', function(actor,event) {
actor.place.launch();
});
}
Main.panel._centerBox.add(this.actor, { y_fill: true });
Main.panel._menus.addMenu(this.menu);
}
};
function main(extensionMeta) {
let userExtensionLocalePath = extensionMeta.path + '/locale';
Gettext.bindtextdomain("places_button", userExtensionLocalePath);
Gettext.textdomain("places_button");
new PlacesButton();
}
The heavy lifting in creating icons is done by iconFactory which is a JavaScript callback that creates an icon texture given a size parameter. It is implemented in ../js/ui/placeDisplay.js
iconFactory: function(size) {
let icon = this._mount.get_icon();
return St.TextureCache.get_default().load_gicon(null, icon, size);
},
Example 5:
In this example, I show you how to modify the previous example to display icons followed by labels on each menu option as shown below:
GNOME3 Shell screenshot
Here is the relevant source code:
const St = imports.gi.St;
const Main = imports.ui.main;
const PopupMenu = imports.ui.popupMenu;
const PanelMenu = imports.ui.panelMenu;
const Gettext = imports.gettext;
const _ = Gettext.gettext;
const MYPLACES_ICON_SIZE = 22;
function MyPopupMenuItem() {
this._init.apply(this, arguments);
}
MyPopupMenuItem.prototype = {
__proto__: PopupMenu.PopupBaseMenuItem.prototype,
_init: function(icon, text, params) {
PopupMenu.PopupBaseMenuItem.prototype._init.call(this, params);
this.addActor(icon);
this.label = new St.Label({ text: text });
this.addActor(this.label);
}
};
function PlacesButton() {
this._init();
}
function PlacesButton() {
this._init();
}
PlacesButton.prototype = {
__proto__: PanelMenu.Button.prototype,
_init: function() {
PanelMenu.Button.prototype._init.call(this, 0.0);
this._icon = new St.Icon({ icon_name: 'start-here',
icon_type: St.IconType.SYMBOLIC,
style_class: 'system-status-icon' });
this.actor.set_child(this._icon);
Main.panel._centerBox.add(this.actor, { y_fill: true });
let placeid;
this.placeItems = [];
this.defaultPlaces = Main.placesManager.getDefaultPlaces();
this.bookmarks = Main.placesManager.getBookmarks();
this.mounts = Main.placesManager.getMounts();
// Display default places
for ( placeid = 0; placeid < this.defaultPlaces.length; placeid++) {
let icon = this.defaultPlaces[placeid].iconFactory(MYPLACES_ICON_SIZE);
this.placeItems[placeid] = new MyPopupMenuItem(icon, _(this.defaultPlaces[placeid].name));
this.menu.addMenuItem(this.placeItems[placeid]);
this.placeItems[placeid].connect('activate', function(actor,event) {
actor.place.launch();
});
}
this.menu.addMenuItem(new PopupMenu.PopupSeparatorMenuItem());
// Display default bookmarks
for ( let bookmarkid = 0; bookmarkid < this.bookmarks.length; bookmarkid++, placeid++) {
let icon = this.bookmarks[bookmarkid].iconFactory(MYPLACES_ICON_SIZE);
this.placeItems[placeid] = new MyPopupMenuItem(icon, _(this.bookmarks[bookmarkid].name));
this.menu.addMenuItem(this.placeItems[placeid]);
this.placeItems[placeid].connect('activate', function(actor,event) {
actor.place.launch();
});
}
if (this.mounts.length > 0) {
this.menu.addMenuItem(new PopupMenu.PopupSeparatorMenuItem());
}
// Display default mounts
for ( let mountid = 0; mountid < this.mounts.length; placeid++, mountid++ ) {
let icon = this.mounts[mountid].iconFactory(MYPLACES_ICON_SIZE);
this.placeItems[placeid] = new MyPopupMenuItem(icon, _(this.mounts[mountid].name) );
this.placeItems[placeid].place = this.mounts[mountid];
this.menu.addMenuItem(this.placeItems[placeid]);
this.placeItems[placeid].connect('activate', function(actor,event) {
actor.place.launch();
});
}
Main.panel._centerBox.add(this.actor, { y_fill: true });
Main.panel._menus.addMenu(this.menu);
}
};
function main(extensionMeta) {
let userExtensionLocalePath = extensionMeta.path + '/locale';
Gettext.bindtextdomain("places_button", userExtensionLocalePath);
Gettext.textdomain("places_button");
new PlacesButton();
}
I had to implement my own version of PopupMenuItem called MyPopupMenuItem in order to display an icon in front on the menuitem label. This is basically just a wrapper around PopupBaseMenuItem.
Example 6:
In this example, I show you how to add an Applications menu next to the Activities button.
GNOME3 Shell screenshot
Here is the source code for extension.js:
const St = imports.gi.St;
const Main = imports.ui.main;
const PopupMenu = imports.ui.popupMenu;
const PanelMenu = imports.ui.panelMenu;
const Shell = imports.gi.Shell;
const Lang = imports.lang;
const Gettext = imports.gettext;
const _ = Gettext.gettext;
const APPMENU_ICON_SIZE = 22;
function MyPopupMenuItem() {
this._init.apply(this, arguments);
}
MyPopupMenuItem.prototype = {
__proto__: PopupMenu.PopupBaseMenuItem.prototype,
_init: function(icon, text, menu_icon_first, params) {
PopupMenu.PopupBaseMenuItem.prototype._init.call(this, params);
this.label = new St.Label({ text: text });
if (menu_icon_first) {
this.box = new St.BoxLayout({ style_class: 'applications-menu-box'});
this.box.add(icon);
this.box.add(this.label);
this.addActor(this.box);
} else {
this.addActor(this.label);
this.addActor(icon);
}
}
};
function ApplicationsButton() {
this._init.apply(this, arguments);
}
ApplicationsButton.prototype = {
__proto__: PanelMenu.Button.prototype,
_init: function(mode) {
PanelMenu.Button.prototype._init.call(this, 0.0);
this._icon = new St.Icon({ icon_name: 'fedora-logo-icon',
icon_type: St.IconType.FULLCOLOR,
icon_size: Main.panel.button.height });
this.actor.set_child(this._icon);
this._appSys = Shell.AppSystem.get_default();
this._categories = this._appSys.get_sections();
this._menuIconFirst = mode;
this._display();
this._appSys.connect('installed-changed', Lang.bind(this, function() {
Main.queueDeferredWork(this._reDisplay);
}));
// add immediately after hotspot
Main.panel._leftBox.insert_actor(this.actor, 1);
Main.panel._menus.addMenu(this.menu);
},
_display: function() {
this.appItems = [];
for (let id = 0; id < this._categories.length; id++) {
this.appItems[this._categories[id]] = new PopupMenu.PopupSubMenuMenuItem(this._categories[id]);
this.menu.addMenuItem(this.appItems[this._categories[id]]);
}
let appInfos = this._appSys.get_flattened_apps().filter(function(app) {
return !app.get_is_nodisplay();
});
for (let appid = appInfos.length-1; appid >= 0; appid--) {
let appInfo = appInfos[appid];
let sec = appInfo.get_section();
if (sec != null) {
let icon = appInfo.create_icon_texture(APPMENU_ICON_SIZE);
let appName = new MyPopupMenuItem(icon, appInfo.get_name(), this._menuIconFirst);
// let appName = new PopupMenu.PopupMenuItem(appInfo.get_name());
appName._appInfo = appInfo;
this.appItems[appInfo.get_section()].menu.addMenuItem(appName);
appName.connect('activate', function(actor,event) {
let id = actor._appInfo.get_id();
Shell.AppSystem.get_default().get_app(id).activate(-1);
});
}
}
},
_redisplay: function() {
for (let id = 0; id < this._categories.length; id++) {
this.appItems[this._categories[id]].menu.destroy();
}
this._display();
}
};
function main(extensionMeta) {
let userExtensionLocalePath = extensionMeta.path + '/locale';
Gettext.bindtextdomain("applications_button", userExtensionLocalePath);
Gettext.textdomain("applications_button");
new ApplicationsButton(false);
}
Again I implemented my own version of PopupMenuItem called MyPopupMenuItem in order to handle displaying an icon in front on the menuitem label or visa versa as shown below. If false is passed to ApplicationsButton when creating the new object, menu options are displayed label followed by icon, otherwise they are displayed icon followed by label.
GNOME3 Shell screenshot
Example 7:
Not all components of the GNOME Shell can be easily modified or customized. For example, suppose I would like to display a search provider’s logo as an icon on their search provider button. The icon is already available as a base64 string in a search provider’s OpenSearch xml file but is not currently used by the GNOME Shell.
Here is how the search provider buttons look like at present:
GNOME3 Shell screenshot
Looking at the current source code for the GNOME Shell, we can see that in search.js the icon string is read by parse_search_provider and stored in the _providers array.
_addProvider: function(fileName) {
let path = global.datadir + '/search_providers/' + fileName;
let source = Shell.get_file_contents_utf8_sync(path);
let [success, name, url, langs, icon_uri] = global.parse_search_provider(source);
let provider ={ name: name,
url: url,
id: this._providers.length,
icon_uri: icon_uri,
langs: langs };
if (this._checkSupportedProviderLanguage(provider)) {
this._providers.push(provider);
this.emit('changed');
}
},
However it is not passed to where the search provider button is created in searchDisplay.js due to a limitation in the following piece of code in search.js:
getProviders: function() {
let res = [];
for (let i = 0; i < this._providers.length; i++)
res.push({ id: i, name: this._providers[i].name });
return res;
},
This needs to be modified to include the icon_uri base64 string, i.e.
getProviders: function() {
let res = [];
for (let i = 0; i < this._providers.length; i++)
res.push({ id: i, name: this._providers[i].name, icon_uri: this._providers[i].icon_uri });
return res;
},
The next problem to overcome is how to decode the icon_uri base64 string and convert it into an actor.to displaying a search provider’s icon on their button. To do this you need to modify _createOpenSearchProviderButton in searchDisplay.js to display the icon as well as the name of the search provider.
Here is the original code
_createOpenSearchProviderButton: function(provider) {
let button = new St.Button({ style_class: 'dash-search-button',
reactive: true,
x_fill: true,
y_align: St.Align.MIDDLE });
let bin = new St.Bin({ x_fill: false,
x_align:St.Align.MIDDLE });
button.connect('clicked', Lang.bind(this, function() {
this._openSearchSystem.activateResult(provider.id);
}));
let title = new St.Label({ text: provider.name,
style_class: 'dash-search-button-label' });
bin.set_child(title);
button.set_child(bin);
provider.actor = button;
this._searchProvidersBox.add(button);
},
and here is the modified code which displays the search providers icon as well as their name.
_createOpenSearchProviderButton: function(provider) {
let button = new St.Button({ style_class: 'dash-search-button',
reactive: true,
x_fill: true,
y_align: St.Align.MIDDLE });
button.connect('clicked', Lang.bind(this, function() {
this._openSearchSystem.activateResult(provider.id);
}));
let title = new St.Label({ text: provider.name,
style_class: 'dash-search-button-label' });
let textureCache = St.TextureCache.get_default();
let searchIcon = textureCache.load_uri_sync(ST_TEXTURE_CACHE_POLICY_FOREVER,
provider.icon_uri, -1, -1);
let iconBin = new St.Bin({ style_class: 'dash-search-button-icon',
child: searchIcon });
let box = new St.BoxLayout();
box.add(iconBin, {expand: true, x_fill: false, x_align: St.Align.END });
box.add(title, {expand: true, x_fill: false, x_align: St.Align.START });
button.set_child(box);
provider.actor = button;
this._searchProvidersBox.add(button);
},
With this modified code, here is how the search provider buttons now look like:
GNOME3 Shell screenshot
A GNOME Shell extension could probably be written to monkey patch the modified versions of the two functions into the GNOME Shell. It is something I will try to do when I get some free time.
Well, this post is getting too big and so it is time to conclude it. But before I do, I want to mention something about the GNOME Shell that has been on my mind recently as I experiment with its internals. For some reason the GNOME developers seem to think that the GNOME Shell should be an integral part of the OS and that is obvious in some of the design decisions. Recently Colin Walters stated that
Where we want to get to is that there are really just three things:
* Apps
* Extensions
* The OS
This is just plain wrong as far as I am concerned. A user should always have a choice of desktop managers and shells. Sounds like the vision of the GNOME developers, at least as articulated by Colin Walters in his post, is that the OS and the GNOME Shell should be one and the same as in Microsoft Windows. If this is the goal, then I fear that many existing users will abandon the GNOME desktop.
Sometimes I get the impression that the GNOME Shell was designed and put together by a group of arrogant young software engineers who were more interested in adhering to the tenants of so-called usability and design theory studies than in what their end users really needed in a modern graphical shell. Frankly, I fully agree with what Bruce Byfield recently wrote in Datamation about usability studies hurting the free desktop.
Fortunately, the decision to use JavaScript for the GNOME Shell UI may prove to be the salvation of the GNOME Shell as it allows ordinary users to quickly experiment with modifying parts of the UI. From what I see on the Internet, there are quite a few people experimenting with customizing the GNOME Shell either by directly modifying the source code or by developing extensions. While GNOME developers such as Lennart Poettering are advocating stricter vertical integration of the platform, I think that the decision to use JavaScript will be key to enabling users of GNOME Shell to easily bypass the diktats of the anointed few.
Don’t get me wrong – I like and use the GNOME Shell and really do not want to go back to GNOME 2 or another desktop manager, but I am often frustrated by it’s design constraints which get in the way of me doing what I want to do quickly and efficiently. Worse, when you look under the hood of the GNOME Shell, you quickly become aware of serious shortcomings in much of the underlying codebase. For example, why is the GNOME shell dependant on two distinct JavaScript engines, i.e. gjs and seed (webkit)? Pick one of these two JavaScript engines and remove the dependencies on the other. And why am I forced to use Evolution for calendaring? What happened to choice?
Adieu, and keep experimenting!
P.S. All of the extensions in this post will be downloadable from here.
88 comments to More GNOME Shell Customization
• Hi,
I tried to modify the ‘autohidetopbar’ extension, because I only want to see the top panel in the Overview screen. I think I only have to edit the lines ‘Main.panel.actor.connect(‘leave-event’, Lang.bind(Main.panel, Main.panel._hidePanel));’ and ‘Main.panel.actor.connect(‘enter-event’, Lang.bind(Main.panel, Main.panel._showPanel));’. The problem is that I don’t know where to connect functions to.
I tried it with things like ‘Main.overview.connect(“show”, Lang.bind(Main.panel, Main.panel._showPanel));’, but that didn’t work.
How can I find out where I can connect those callbacks to?
• Rom
all I can say : THANKS !
• puya
Hi
Thanks for great examples,
just a question, How can i write unicode characters on the top panel?
Thanks in advance!!
• Should be the same as you would do to write a UTF-8 string to the top panel. I have not tried it but if you provide an example of what you are trying to do, I will try and assist you.
• Stefan
Hi,
nice tutorials. I have a question too:
Do you know how to get subfolders or files of a folder (e.g. home). I don’t find anything in lg and dunno where to find an API or something like that.
As programmer it’s easy to follow your examples, but hard to write something new, cause there is no docs or is there?
• mnv
Thanks for good examples!
Where may I read documentation for used objects in this examples?
const St = imports.gi.St;
const Main = imports.ui.main;
const PopupMenu = imports.ui.popupMenu;
const PanelMenu = imports.ui.panelMenu;
const Gettext = imports.gettext;
• Thanks for this: I was able to modify it to add my favourite apps’ systray icons to the panel.
Some comments:
1. You declare const Panel, but don’t use it; why?
2. The standard format for extensions now seems to be to have init, enable and disable methods. I tried renaming your main method to enable, providing a blank init method, but it doesn’t work; also, I tried to add a disable method which set the various elements to null instead of a string, but that didn’t work either. Any hints on making this extension better-behaved (so it can e.g. be turned on and off in gnome-tweak-tool?).
3. The comment says “Use the english text string displayed when hovering your mouse over the bottom right notification area”, but this seems to be wrong: for example, for caffeine (the screensaver controller), the text string is caffeine-cup-empty (even when it’s full!) but the string to use is “caffeine”. Similarly, the string displayed for GMPC is “GNOME Music Player” or something, but the string to use is “gmpc”. Hence, the correct string would seem to be the application name; at least, that’s more likely to work.
4. Could you please put a mailto: link on your contact page? Using web forms is a pain as I don’t get to keep a record of what I wrote. (The way to avoid spam is to use a spam filter, not to hide your email address!)
I tried sending this as an email, but I got an error “Illegal characters in POST. Possible email injection attempt.”
• Preston
First and foremost your Gnome 3 extensions are brilliant and much appreciated! Thanks for the tweaks and tutorial.
I switched to Gnome 3 on Ubuntu 11.10 and hiding the top bar is the feature I desire Most! I installed your “autohidetopbar” extension and it didn’t work at first. Using looking glass I got errors for missing init, enable, and disable methods (is that new on the Gnome 3 moving target, or just an Ubuntu distro issue…). “As a test” I modified the extension to add empty “enable” and “disable” methods, then created “init” to call “main”. I also updated the shell version in .json to 3.2.0.
That made the Top bar hide on double click, but then I couldn’t get it to reveal again. So I was curious, do you have a version of the autohidetopbar extension that addresses the issue I’m encountering?
Thanks again!
• I’m sure you’re sick of requests for this by now, but could you please update the autohide top bar extension for Gnome 3.2 in Ubuntu 11.10? I can confirm that updating the version in .json doesn’t help. I even tried moving the autohide extension to /usr/share/gnome-shell/extensions with all the other gnome-shell extensions. It shows up in the gnome-tweak shell extension section, but it has an exclamation mark next to it and the on/off switch is grayed out.
Love your site and all your hard work on Gnome-shell, otherwise. You single handedly made Gnome Shell usable for me.
• Bazzer
As above, another vote for some work to the autohidetopbar in Ubuntu. Thanks for explaining all JS above, too.
• Najam
Any help to update the battery status icon on leftbox to show remaining battery life ( % or hr:min ) adjacent to it ? Now I need to click it everytime to see it. I tried my luck, in vain :)
BTW, your blog is just great !!
• wnfu12nb
Hi,
I tried to modify the ‘autohidetopbar’ extension, because I only want to see the top panel in the Overview screen. I think I only have to edit the lines ‘Main.panel.actor.connect(‘leave-event’, Lang.bind(Main.panel, Main.panel._hidePanel));’ and ‘Main.panel.actor.connect(‘enter-event’, Lang.bind(Main.panel, Main.panel._showPanel));’. The problem is that I don’t know where to connect functions to.
I tried it with things like ‘Main.overview.connect(“show”, Lang.bind(Main.panel, Main.panel._showPanel));’, but that didn’t work.
How can I find out where I can connect those callbacks to?
• You need to read the source code under /usr/share/gnome-shell/js/ui
• ykkhern
Hi,
I just modified “extension.js” as per following:
1. Line 104, change “main” to “init” –> function init() {
2. add the following at the end of file:
function enable() {
}
function disable() {
}
Then modify Line 2 of “metadata.json” to:
“shell-version”: ["3.2"]
Well, it’s kind of working although the top panel is not really hideaway that only the background is transparent and the rest is in very light gray. All windows can be maximized to the top of screen and you still can see topbar in Overview screen.
If you move the mouse cursor all the way to the top of screen (without going into overview screen) you still can use any item on the topbar like checking calender, battery status, changing volume, and etc.
It is not perfect yet but for a person like me who don’t know a thing about java script, this modified extension is now working great.
• Bastiaan
Hi,
Great info!
I have 2 whiches you might be able to help me:
1) Is there a way to add the ubuntu software center (or any software repository for that matter) to the search result in the activity window?
2) How do I start an application ion the top right hot spot? Like the circular-application-menu (http://code.google.com/p/circular-application-menu/)
Thanks!
Bastiaan
• bmbaker
Hi there :-)
i was wondering if you will be updating your theme extension for 3.2 ?
cheers :-)
BB
• Really nice,
very cool sharing.
|
__label__pos
| 0.93259 |
I'm a Civil Engineer in Computation and Computer Science of Andrés Bello University, Santiago - Chile. I'm working as a Software Engineer in Ingerencia, an small company that develop software solutions. In paralell i'm working on a project with a new company that calls ID (Ingeniería Domótica, in english means Home Automation Engineering), where i'm developing a software solution based on Netbeans platform 6.8 and some other interesting technologies like PLCBUS protocol and hardware devices. I love Linux and I love Java. In my other life I'm a heavy weight lifter, 7 years already. What can I say, I'm combining the develop of the mind and body. Daniel has posted 3 posts at DZone. View Full User Profile
How to Create a Custom NetBeans Platform LifecycleManager
06.07.2010
| 4874 views |
• submit to reddit
One of the things we need when we develop a desktop application is closing control: "What should our application do if we press the exit button? Should we save our work? Should we confirm the exit? Or should we just close the application and lose everything?"
If you are using the NetBeans Platform, it is possible that you have found yourself in this dilemma. The solution is VERY simple.
First, we need to create a class with an annotation for overriding the NetBeans Platform's default org.openide.LifecycleManager. Call the class "MyLifeCycleManager" and then annotate the class as follows:
@ServiceProvider(service=LifecycleManager.class, position=1)
Then we need to add the code for our class, for example:
@ServiceProvider(service=LifecycleManager.class, position=1)
public class MyLifecycleManager extends LifecycleManager {
@Override
public void saveAll() {
}
@Override
public void exit() {
Set<Topcomponent> tcs = TopComponent.getRegistry().getOpened();
Iterator<Topcomponent> it = tcs.iterator();
while (it.hasNext()) {
TopComponent tc = it.next();
if (tc instanceof DepartmentEditorTopComponent) {
tc.close();
}
}
Collection c = Lookup.getDefault().lookup(new Lookup.Template(LifecycleManager.class)).allInstances();
for (Iterator i = c.iterator(); i.hasNext();) {
LifecycleManager lm = (LifecycleManager) i.next();
if (lm != this) {
lm.exit();
}
}
}
}
And this is how it looks in the IDE:
When we extend "LifecycleManager", we need to implement the "exit" method and it is in that implementation where we code what we want to do before closing our application. In this case, I worry about closing properly a TopComponent I made (DepartmentTopComponent, from my Home Automation Application). You MUST write the following piece of code in the last part of exit method, so that your application can close properly (System.exit(0) will not work):
Collection c = Lookup.getDefault().lookup (new Lookup.Template (LifecycleManager.class)).allInstances();
for (Iterator i=c.iterator(); i.hasNext();) {
LifecycleManager lm = (LifecycleManager) i.next();
if (lm != this) {
lm.exit();
}
}
With that, you have your own implementation of LifecycleManager and you can provide specific things that need to be done before closing your application.
0
Your rating: None
Published at DZone with permission of its author, Daniel Morales.
(Note: Opinions expressed in this article and its replies are the opinions of their respective authors and not those of DZone, Inc.)
Comments
Rupert Young replied on Sat, 2014/07/05 - 4:38pm
Hello,
This looks like what I need but when I click exit I get this error:
"java.lang.AssertionError: Window System API is required to be called from AWT thread"
Any ideas for a resolution?
Regards,
Rupert
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Components
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Rendering-iterables
You will often want to display multiple similar components from a collection of data. The rx.foreach component takes an iterable (list, tuple or dict) and a function that renders each item in the list. This is useful for dynamically rendering a list of items defined in a state.
In this first simple example we iterate through a list of colors and render the name of the color and use this color as the background for that rx.box. As we can see we have a function colored_box that we pass to the rx.foreach component. This function renders each item from the list that we have defined as a state var color.
red
green
blue
yellow
orange
purple
class IterState(rx.State):
color: list[str] = [
"red",
"green",
"blue",
"yellow",
"orange",
"purple",
]
def colored_box(color: str):
return rx.box(rx.text(color), background_color=color)
def simple_foreach():
return rx.grid(
rx.foreach(IterState.color, colored_box),
columns="6",
)
The function can also take an index as a second argument, meaning that we can enumerate through data as shown in the example below.
0
1
2
3
4
5
class IterIndexState(rx.State):
color: list[str] = [
"red",
"green",
"blue",
"yellow",
"orange",
"purple",
]
def enumerate_foreach():
return rx.grid(
rx.foreach(
IterIndexState.color,
lambda color, index: rx.box(
rx.text(index), bg=color
),
),
columns="6",
)
We can iterate through a dict data structure using a foreach. When the dict is passed through to the function that renders each item, it is presented as a list of key-value pairs [("sky", "blue"), ("balloon", "red"), ("grass", "green")].
sky
balloon
grass
class SimpleDictIterState(rx.State):
color_chart: dict[str, str] = {
"sky": "blue",
"balloon": "red",
"grass": "green",
}
def display_color(color: list):
# color is presented as a list key-value pairs [("sky", "blue"), ("balloon", "red"), ("grass", "green")]
return rx.box(
rx.text(color[0]), bg=color[1], padding_x="1.5em"
)
def dict_foreach():
return rx.grid(
rx.foreach(
SimpleDictIterState.color_chart,
display_color,
),
columns="3",
)
rx.foreach can be used with nested state vars. Here we use nested foreach components to render the nested state vars. The rx.foreach(project["technologies"], get_badge) inside of the project_item function, renders the dict values which are of type list. The rx.box(rx.foreach(NestedStateFE.projects, project_item)) inside of the projects_example function renders each dict inside of the overall state var projects.
Next.jsPrismaTailwindGoogle CloudDockerMySQL
PythonFlaskGoogle CloudDocker
class NestedStateFE(rx.State):
projects: list[dict[str, list]] = [
{
"technologies": [
"Next.js",
"Prisma",
"Tailwind",
"Google Cloud",
"Docker",
"MySQL",
]
},
{
"technologies": [
"Python",
"Flask",
"Google Cloud",
"Docker",
]
},
]
def get_badge(technology: str) -> rx.Component:
return rx.badge(
technology, variant="soft", color_scheme="green"
)
def project_item(project: dict) -> rx.Component:
return rx.box(
rx.hstack(
rx.foreach(project["technologies"], get_badge)
),
)
def projects_example() -> rx.Component:
return rx.box(
rx.foreach(NestedStateFE.projects, project_item)
)
If you want an example where not all of the values in the dict are the same type then check out the example on var operations using foreach .
Here is a further example of how to use foreach with a nested data structure.
purple
red
blue
orange
yellow
red
green
blue
yellow
class NestedDictIterState(rx.State):
color_chart: dict[str, list[str]] = {
"purple": ["red", "blue"],
"orange": ["yellow", "red"],
"green": ["blue", "yellow"],
}
def display_colors(color: list[str, list[str]]):
return rx.vstack(
rx.text(color[0], color=color[0]),
rx.hstack(
rx.foreach(
color[1],
lambda x: rx.box(
rx.text(x, color="black"), bg=x
),
)
),
)
def nested_dict_foreach():
return rx.grid(
rx.foreach(
NestedDictIterState.color_chart,
display_colors,
),
columns="3",
)
We can also use foreach with the cond component.
In this example we define the function render_item. This function takes in an item, uses the cond to check if the item is_packed. If it is packed it returns the item_name with a next to it, and if not then it just returns the item_name. We use the foreach to iterate over all of the items in the to_do_list using the render_item function.
Sammy's Packing List
• Space suit ✔
• Helmet ✔
• Back Pack
class ToDoListItem(rx.Base):
item_name: str
is_packed: bool
class ForeachCondState(rx.State):
to_do_list: list[ToDoListItem] = [
ToDoListItem(
item_name="Space suit", is_packed=True
),
ToDoListItem(item_name="Helmet", is_packed=True),
ToDoListItem(
item_name="Back Pack", is_packed=False
),
]
def render_item(item: [str, bool]):
return rx.cond(
item.is_packed,
rx.list.item(item.item_name + " ✔"),
rx.list.item(item.item_name),
)
def packing_list():
return rx.vstack(
rx.text("Sammy's Packing List"),
rx.list(
rx.foreach(
ForeachCondState.to_do_list, render_item
)
),
)
← Conditional RenderingLibrary →
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Atomic Nucleus
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Atomic Nucleus
An atomic bomb explodes and generates a huge mushroom cloud. The tremendous energy released when the bomb explodes is incredibly destructive. Where does all the energy come from? The answer is the nucleus of the atom.
At the Heart of It All
The nucleus (plural, nuclei) is a positively charged region at the center of the atom. It consists of two types of subatomic particles packed tightly together. The particles are protons, which have a positive electric charge, and neutrons, which are neutral in electric charge. Outside of the nucleus, an atom is mostly empty space, with orbiting negative particles called electrons whizzing through it. The Figure below shows these parts of the atom.
Size and Mass of the Nucleus
The nucleus of the atom is extremely small. Its radius is only about 1/100,000 of the total radius of the atom. If an atom were the size of a football stadium, the nucleus would be about the size of a pea! You can see another sports analogy comparing the size of the nucleus and atom at this URL:
http://www.youtube.com/watch?v=1o-FsxAkvZk (1:06)
Electrons have virtually no mass, but protons and neutrons have a lot of mass for their size. As a result, the nucleus has virtually all the mass of an atom. Given its great mass and tiny size, the nucleus is very dense. If an object the size of a penny had the same density as the nucleus of an atom, its mass would be greater than 30 million tons! You can learn more about the size and mass of the nucleus at this URL:
http://www.youtube.com/watch?v=Tfy0sIVfVOY (2:03)
Holding It All Together
Particles with opposite electric charges attract each other. This explains why negative electrons orbit the positive nucleus. Particles with the same electric charge repel each other. This means that the positive protons in the nucleus push apart from one another. So why doesn’t the nucleus fly apart? An even stronger force—called the strong nuclear force—holds protons and neutrons together in the nucleus. You can learn more about these forces in the nucleus by watching the video at this URL:
http://www.youtube.com/watch?v=PdFsb2sWrW4 (6:27)
Q: Can you guess why an atomic bomb releases so much energy when it explodes?
A: When an atomic bomb explodes, the nuclei of atoms undergo a process called fission, in which they split apart. This releases the huge amount of energy that was holding together subatomic particles in the nucleus.
Summary
• The nucleus is a small, dense region at the center of the atom. It consists of positive protons and neutral neutrons, so it has an overall positive charge.
• The nucleus is just a tiny part of the atom, but it contains virtually all of the atom’s mass.
• The strong nuclear force holds together protons and neutrons in the nucleus and overcomes the electric force of repulsion between protons.
Vocabulary
• nucleus : Tiny region in the center of an atom that contains protons and neutrons and makes up almost all of the atom’s mass.
Practice
Watch this short video about how the nucleus was discovered, and then answer the questions below.
http://www.youtube.com/watch?v=Q8RuO2ekNGw (0:48)
1. Describe the scientific procedure that was used to discover the nucleus.
2. What evidence led scientists to conclude that atoms consist mostly of empty space with a very small, positively charged mass at the center?
3. Reflect on the method used in the experiment. Why was it important to send positive—as opposed to neutral or negative—particles toward the gold foil?
Review
1. Describe the nucleus of the atom.
2. Why is the nucleus positive in charge?
3. Explain why the nucleus is very dense.
4. Outline the forces that act on particles in the nucleus.
5. If you made a three-dimensional model of an atom and its nucleus, how would you represent the atom? How would you represent nucleus? Explain your choices.
Vocabulary
nucleus
nucleus
Tiny region in the center of an atom that contains protons and neutrons and makes up almost all of the atom’s mass.
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