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--- |
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license: openrail++ |
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base_model: stabilityai/stable-diffusion-xl-base-1.0 |
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instance_prompt: a picture of <s1><s2> minifigure |
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tags: |
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- stable-diffusion-xl |
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- stable-diffusion-xl-diffusers |
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- text-to-image |
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- diffusers |
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- lora |
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inference: true |
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--- |
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# LoRA DreamBooth - merve/lego-lora-trained-xl |
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These are LoRA adaption weights for stabilityai/stable-diffusion-xl-base-1.0. The weights were trained on a picture of <s1><s2> minifigure using [DreamBooth](https://dreambooth.github.io/). You can find some example images in the following. |
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![img_0](./image_0.png) |
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![img_1](./image_1.png) |
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![img_2](./image_2.png) |
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![img_3](./image_3.png) |
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You can use this code ๐ |
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```python |
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from huggingface_hub.repocard import RepoCard |
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from diffusers import DiffusionPipeline |
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import torch |
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lora_model_id = "merve/lego-lora-trained-xl" |
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card = RepoCard.load(lora_model_id) |
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base_model_id = card.data.to_dict()["base_model"] |
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pipe = DiffusionPipeline.from_pretrained(base_model_id, torch_dtype=torch.float16) |
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pipe = pipe.to("cuda") |
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pipe.load_lora_weights(lora_model_id) |
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pipe("a picture of <s1><s2> minifigure as lana del rey, high quality", num_inference_steps=35).images[0] |
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``` |
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LoRA for the text encoder was enabled: False. |
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Special VAE used for training: madebyollin/sdxl-vae-fp16-fix. |
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